http://wiki.math.uwaterloo.ca/statwiki/api.php?action=feedcontributions&user=Nkadiyar&feedformat=atomstatwiki - User contributions [US]2022-08-07T22:39:44ZUser contributionsMediaWiki 1.28.3http://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/_Automated_Curriculum_Learning_for_Neural_Networks&diff=31621STAT946F17/ Automated Curriculum Learning for Neural Networks2017-11-28T18:08:59Z<p>Nkadiyar: </p>
<hr />
<div>= Introduction =<br />
<br />
Humans and animals learn much better when the examples are not randomly presented but organized in a meaningful order which illustrates gradually more concepts, and gradually more complex ones. Here, we formalize such training strategies in the context of machine learning, and call them “curriculum learning”. The idea of training a learning machine with a curriculum can be traced back<br />
at least to Elman (1993). The basic idea is to start small, learn easier aspects of the task or easier sub-tasks, and then gradually increase the difficulty level. <br />
<br />
However curriculum learning has only recently become prevalent in the field (e.g., Bengio et al., 2009), due in part to the greater complexity of problems now being considered. In particular,<br />
recent work on learning programs with neural networks has relied on curricula to scale up to longer or more complicated tasks (Reed and de Freitas, 2015, Gui et al. 2017). We expect this trend<br />
to continue as the scope of neural networks widens, with deep reinforcement learning providing fertile ground for structured learning.<br />
<br />
One reason for the slow adoption of curriculum learning is that it’s effectiveness is highly sensitive to the mode of progression through the tasks. One popular approach is to define a hand-chosen performance threshold for advancement to the next task, along with a fixed probability of returning to earlier tasks, to prevent forgetting (Sutskever and Zaremba, 2014). However, as well as introducing hard-to-tune parameters, this poses problems for curricula where appropriate thresholds may be unknown or variable across tasks.<br />
<br />
The main contribution of the paper is that a stochastic policy, continuously adapted to optimize learning progress is proposed. Given a progress signal that can be evaluated for each<br />
training example, we use a multi-armed bandit algorithm to find a stochastic policy over the tasks that maximizes overall progress. The bandit is non-stationary because the behaviour of the network, and hence the optimal policy, evolves during training. Moreover variants of prediction gain, and also a novel class of progress signals which we refer to as complexity gain are considered in this paper.<br />
<br />
= Model =<br />
A task is a distribution $D$ over sequences from $\mathcal{X}$ . A curriculum is an ensemble of tasks $D_1, \ldots D_N$ , a sample is an example drawn from one of the tasks of the curriculum,<br />
and a syllabus is a time-varying sequence of distributions over tasks. A neural network is considered as a probabilistic model $p_\theta$ over $\mathcal{X}$, whose parameters are denoted $\theta$.<br />
<br />
The expected loss of the network on the $k$-th task is <br />
\[<br />
\mathcal{L}_k( \theta) := \mathbb{E}_{\mathbf{x} \sim D_k} L(\mathbf{x}, \theta),<br />
\]<br />
where $L(\mathbf{x}, \theta):= -\log_{p_\theta}(\mathbf{x})$ is the sample loss on $\mathbf{x}$. <br />
<br />
A curriculum containing $N$ tasks as an $N$-armed bandit, and a syllabus as an adaptive policy which seeks to maximize payoffs from this bandit. In the bandit setting, an agent selects a sequence of arms (actions) $a_1,\ldots, a_T$ over T rounds of play. After each round, the selected arm yields a payoff $r_t$; the payoffs for the other arms are not observed.<br />
===Adversarial Multi-Armed Bandits===<br />
The classic algorithm for adversarial bandits is Exp3, which minimize regret with respect to the single best arm evaluated over the whole history. However, in the case of training neural network, an arm is optimal for a portion of the history, then another arm, and so on; the best strategy is then piecewise stationary. The Fixed Share method addresses this issue by using an $\epsilon$-greedy strategy. It is known as the Exp3.S algorithm. <br />
<br />
On round $t$, the agent selects an arm stochastically according to a policy $\pi_t$ . This policy is defined by a set of weights $w_t$,<br />
\[<br />
\pi_t(i) := (1-\epsilon)\frac{e^{w_{t,i}}}{\sum_{j=1}^N e^{w_{t,j}}}+\frac{\epsilon}{N} <br />
\]<br />
\[<br />
w_{t,i}:= \log \big[ (1-\alpha_t)\exp\{ w_{t-1,i} +\eta \bar{r}_{t-1,i}^\beta \} +\frac{\alpha_t}{N-1}\sum_{j \ne i} \exp\{ w_{t-1,j} +\eta \bar{r}_{t-1,j}^\beta \} \big]<br />
\]<br />
\[<br />
w_{1,i} = 0, \quad \alpha_t = t^{-1} , \quad \bar{r}_{s,i}^\beta = \frac{r_s \mathbb{I}_{[a_s = i]}+ \beta}{ \pi_s(i) }<br />
\]<br />
<br />
===Reward Scaling===<br />
The appropriate step size $\eta$ depends on the magnitudes of the rewards, which may not be known. The magnitude of reward depends strongly on the gain signal used to measure<br />
learning progress, as well as varying over time as the model learns. To address this issue, all rewards are adaptively rescale to $[-1,1]$ by<br />
\[<br />
r_t = \begin{cases}<br />
-1 &\quad \text{if } \hat{r}_t < q^{l}_t\\<br />
1 &\quad \text{if } \hat{r}_t > q^{h}_t\\<br />
\frac{2(\hat{r}_t-q^l_t)}{q^h_t-q^l_t} -1 , &\quad \text{otherwise.}<br />
\end{cases}<br />
\]<br />
where $q^l_t$ and $q^h_t$ are the quantiles of the history of unscaled rewards up to time $t$. The authors chose them to be $20$-th and $80$-th percentiles respectively.<br />
<br />
===Algorithm===<br />
The automated curriculum learning is summarized as followed,<br />
<br />
where $\tau(\mathbf{x})$ is the length of the longest input sequence. Since the processing time of each task may be different.<br />
<center><br />
[[File:alg.png | frame | center |Fig 1. Automated Curriculum Learning Algorithm ]]<br />
</center><br />
<br />
= Learning Progress Signals =<br />
The learning progress is the measurement of effect of a training sample on the target objective. It usually is computationally undesirable or even impossible to obtain. Therefore the authors consider a range of signals derived from two distinct indicators of learning progress: 1) loss-driven, in the sense that they equate progress with a decrease in some loss; or 2) complexity-driven, when they equate progress with an increase in model complexity.<br />
<br />
=== Loss-driven Progress===<br />
The loss-driven progress signals compare the predictions made by the model before and after training on some sample $\mathbf{x}$. <br />
<br />
'''Prediction gain(PG)'''<br />
Prediction gain is defined as the instantaneous change in loss for a sample $\mathbf{x}$<br />
\[<br />
v_{PG}:=L(\mathbf{x},\theta)-L(\mathbf{x},\theta')<br />
\]<br />
<br />
'''Gradient prediction gain (GPG)'''<br />
This measures the magnitude of the gradient vector, which has been used an indicator of salience in the active learning samples<br />
\[<br />
v_{GPG}:= || \triangledown L(\mathbf{x},\theta)||^2_2<br />
\]<br />
<br />
'''Self prediction gain (SPG)'''<br />
Self prediction gain samples a second time from the same task to address the bias problem of PG<br />
\[<br />
v_{SPG}:=L(\mathbf{x}',\theta)-L(\mathbf{x}',\theta'), \quad \mathbf{x}' \sim D_k<br />
\]<br />
<br />
<br />
'''Target prediction gain (TPG)'''<br />
\[<br />
v_{TPG}:=L(\mathbf{x}',\theta)-L(\mathbf{x}',\theta'), \quad \mathbf{x}' \sim D_N<br />
\]<br />
Although this estimator might seem like the most accurate measure so far, it tends to suffer from high variance. <br />
<br />
'''Mean prediction gain (MPG)'''<br />
\[<br />
v_{MPG}:=L(\mathbf{x}',\theta)-L(\mathbf{x}',\theta'), \quad \mathbf{x}' \sim D_k, k \sim U_N,<br />
\]<br />
where $U_N$ denotes the uniform distribution on $\{1,\ldots,N\}$.<br />
<br />
===Complexity-driven Progress===<br />
The intuition for complexity gains are derived from the Minimum Description Length principle (Grunwald, 2007). It states that in order to achieve high generalization from a particular dataset, both the number of bits required to describe the model parameters and the data. This is of significance only if the increase in model complexity strongly reduces the data cost. In the case of neural networks, MDL training is done via stochastic variational inference ( Hinton and Van Camp 1993; Graves, 2011; Kingma et al., 2015; Blundell et al., 2015)<br />
<br />
In the stochastic variational inference framework, a variational posterior $P_\phi(\theta)$ over the network weights is maintained during training, with a single weight sample drawn for each training example. An adaptive prior $Q_\psi(\theta)$ are reused for every network weight. In this paper, both of $P$ and $Q$ are set as a diagonal Gaussian distribution, such the complexity cost can be computed analytically <br />
\[<br />
KL(P_\phi|| Q_\psi) = \frac{(\mu_\phi-\mu_\psi)^2+\sigma^2_\phi-\sigma^2_\psi}{２\sigma^2_\psi}+\ln\Big( \frac{\sigma_\psi}{\sigma_\phi} \Big)<br />
\]<br />
<br />
'''Variational complexity gain (VCG)'''<br />
\[v_{VCG}:= KL(P_{\phi'}|| Q_{\psi'}) - KL(P_\phi|| Q_\psi)\]<br />
<br />
'''Gradient variational complexity gain (GVCG)'''<br />
\[<br />
v_{GVCG}:= [\triangledown_{\phi,\psi} KL(P_\phi|| Q_\psi)]^T \triangledown_\phi \mathbb{E}_{\theta \sim P_\phi} L(\mathbf{x},\theta)<br />
\]<br />
<br />
'''L2 gain (L2G)'''<br />
\[<br />
v_{L2G}:=|| \theta' ||^2_2 -|| \theta ||^2_2, \quad v_{GL2G}:=\theta^T [\triangledown_\theta L(\mathbf{x}, \theta)]<br />
\]<br />
<br />
= Experiments =<br />
To test the proposed approach, the authors applied all the gains to three tasks suites: $n$-gram models, repeat copy, and the bAbI tasks.<br />
<br />
Unidirectional LSTM network architecture was used for all experiments, and cross-entropy was used as loss function. The neural network was optimized by RMSProp with momentum of $0.9$<br />
and a learning rate $10^{-5}$. The parameters for Exp3S algorithm were $\eta = 10^{-3}, \beta = 0, \epsilon = 0.05$. All experiments were repeated $10$ times with different random initializations of network weights. Two performance benchmarks are 1) a fixed uniform policy over all the tasks and 2) directly training on the target task (where applicable).<br />
<br />
===N-Gram Language Modelling===<br />
The first experiment is the character-level KneserNey $n$-gram models (Kneser and Ney, 1995) on the King James Bible data from the Canterbury corpus, with the maximum depth parameter $n$ ranging<br />
between $0$ to $10$. It should be notice that the amount of linguistic structure increases monotonically with $n$. <br />
<br />
Fig. 2 shows that most of the complexity-based gain signals (L2G, GL2G, GVCG) progress rapidly through the curriculum before focusing strongly on the $10$-gram task. The loss-driven progress<br />
(PG, GPG, SPG, TG) also tend to move towards higher $n$, although more slowly and with less certainty.<br />
<br />
<center><br />
[[File:ngram.png | frame | center |Fig 2. N-gram policies for different gain signals, truncated at $2 \times 10^8$ steps. All curves are averages over $10$ runs ]]<br />
</center><br />
<br />
===Repeat Copy===<br />
In the repeat copy task (Graves et al., 2014), the network is asked to repeat a random sequence a given number of times. Fig. 3 shows that GVCG solves the target task about twice<br />
as fast as uniform sampling for VI training, and that the PG, SPG and TPG gains are somewhat faster than uniform, especially in the early stages.<br />
<br />
<center><br />
[[File:rcode.png | frame | center |Fig 3. Target task loss (per output) truncated at $1.1 \times 10^9$ steps.. All curves are averages over $10$ runs ]]<br />
</center><br />
<br />
===BAbI===<br />
The bAbI dataset (Weston et al., 2015) includes $20$ synthetic question-answering problems designed to test the basic reasoning capabilities of machine learning models. BAbI was not specifically designed for curriculum learning, but some of the tasks follow a natural ordering, such as ‘Two Arg Relations’, ‘Three Arg Relations’. The authors hope that an efficient syllabus could be<br />
found for learning the whole set.<br />
<br />
Fig. 4 shows that prediction gain (PG) clearly improved on uniform sampling in terms of both learning speed and number of tasks completed; for SPG the same benefits were visible, though less pronounced. The other gains were either roughly equal to or worse than uniform.<br />
<br />
<center><br />
[[File:babi.png | frame | center |Fig 4. Completion curves for the bAbI curriculum, truncated at $3.5\times10^8$ steps. All curves are averages over $10$ runs ]]<br />
</center><br />
<br />
= Conclusion =<br />
1. The experiments suggest that a stochastic syllabus can improve the significant gains, when a suitable progress signal is used.<br />
<br />
2. The uniform random task is a surprisingly strong benchmark.<br />
<br />
3. The learning progress is best evaluated on a local, rather than global basis. In maximum likelihood training, prediction gain is the most consistent signal, while in variational inference<br />
training gradient variational complexity gain performed best.<br />
<br />
= Critique =<br />
In curriculum learning, a popular approach is to define a hand-chosen performance threshold for advancement to the next task, along with a fixed probability of returning to earlier tasks, to prevent forgetting. It is interesting to compare the perform of it and the proposed automated curriculum learning methods.<br />
<br />
= Source =<br />
Graves, Alex, et al. "Automated Curriculum Learning for Neural Networks." arXiv preprint arXiv:1704.03003 (2017).<br />
<br />
Elman, Jeffrey L. "Learning and development in neural networks: The importance of starting small." Cognition 48.1 (1993): 71-99.<br />
<br />
Bengio, Yoshua, et al. "Curriculum learning." Proceedings of the 26th annual international conference on machine learning. ACM, 2009.<br />
<br />
Reed, Scott, and Nando De Freitas. "Neural programmer-interpreters." arXiv preprint arXiv:1511.06279 (2015).<br />
<br />
Gui, Liangke, Tadas Baltrušaitis, and Louis-Philippe Morency. "Curriculum Learning for Facial Expression Recognition." Automatic Face & Gesture Recognition (FG 2017), 2017 12th IEEE International Conference on. IEEE, 2017.<br />
<br />
Zaremba, Wojciech, and Ilya Sutskever. "Learning to execute." arXiv preprint arXiv:1410.4615 (2014).<br />
<br />
Kneser, Reinhard, and Hermann Ney. "Improved backing-off for m-gram language modeling." Acoustics, Speech, and Signal Processing, 1995. ICASSP-95., 1995 International Conference on. Vol. 1. IEEE, 1995.<br />
<br />
Weston, Jason, et al. "Towards ai-complete question answering: A set of prerequisite toy tasks." arXiv preprint arXiv:1502.05698 (2015).<br />
<br />
Graves, Alex, Greg Wayne, and Ivo Danihelka. "Neural turing machines." arXiv preprint arXiv:1410.5401 (2014).<br />
<br />
Grunwald, P. D. (2007). The minimum description length principle. The MIT Press.<br />
<br />
Hinton, G. E. and Van Camp, D. (1993). Keeping the neural networks simple by minimizing the description length of the weights. In Proceedings of the sixth annual conference on Computational learning theory, pages 5–13. ACM<br />
<br />
Blundell, C., Cornebise, J., Kavukcuoglu, K., and Wierstra, D. (2015). Weight uncertainty in neural networks. In Proceedings of The 32nd International Conference on Machine Learning, pages 1613–1622<br />
<br />
Kingma, D. P., Salimans, T., and Welling, M. (2015). Variational dropout and the local reparameterization trick. In Advances in Neural Information Processing Systems, pages 2575–2583</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/Decoding_with_Value_Networks_for_Neural_Machine_Translation&diff=31619STAT946F17/Decoding with Value Networks for Neural Machine Translation2017-11-28T18:07:07Z<p>Nkadiyar: </p>
<hr />
<div>=Introduction=<br />
<br />
==Background Knowledge==<br />
*'''NTM'''<br />
Neural Machine Translation (NMT), which is based on deep neural networks and provides an end- to-end solution to machine translation, uses an '''RNN-based encoder-decoder architecture''' to model the entire translation process. Specifically, an NMT system first reads the source sentence using an encoder to build a "thought" vector, a sequence of numbers that represents the sentence meaning; a decoder, then, processes the "meaning" vector to emit a translation. (Figure 1)<sup>[[#References|[1]]]</sup><br />
[[File:VNFigure1.png|thumb|600px|center|Figure 1: Encoder-decoder architecture – example of a general approach for NMT.]]<br />
<br />
<br />
<br />
*'''Sequence-to-Sequence(Seq2Seq) Model'''<br />
[[File:VNFigure4.png|thumb|450px|center|Figure 2: Seq2Seq Model]]<br />
- Two RNNs: an encoder RNN, and a decoder RNN<br />
<br />
#The input is passed though the encoder and it’s final hidden state, the “thought vector” is passed to the decoder as it’s initial hidden state. <br />
#Decoder given the start of sequence token, <SOS>, and iteratively produces output until it outputs the end of sequence token, <EOS><br />
<br />
- Commonly used in text generation, machine translation, and related problems<br />
<br />
<br />
<br />
*'''Beam Search'''<br />
Decoding process:<br />
[[File:VNFigure2.png|thumb|600px|center|Figure 3]]<br />
Problem: Choosing the word with highest score at each time step t is not necessarily going to give you the sentence with the highest probability(Figure 3). Beam search solves this problem (Figure 4). Beam search is a heuristic search algorithm such that at each time step t, it takes the top m proposal and continues decoding with each one of them. In the end, you will get a sentence with the highest probability not in the word level. The algorithm terminates if all sentences are completely generated, i.e., all sentences are ended with the <EOS> token.<br />
[[File:VNFigure3.png|thumb|600px|center|Figure 4]]<br />
<br />
==value network==<br />
<br />
Upon the success of NMT with beam search, beam search tends to focus more on short-term reward, which is called myopic bias. At time t, a word $w$ may be appended to the candidates $y_{<t-1} = y_1,...,y_{t}$ if $P(y_{<t}+w|x) > P(y_{<t}+w'|x)$ even if the word $w'$ is the ground truth translation at step t or can offer a better score in future decodings. By applying the concept of value function in Reinforcement Learning (RL), the authors develop a neural network-based prediction model, '''the value network for NMT''', to estimate the long-term reward when appending $w$ to $y_{<t}$ to address the myopic bias.<br />
<br />
The value network takes the source sentence and any partial target sequence as input, and outputs a predicted value to estimate the expected total reward (e.g. BLEU<sup>[[#References|[2]]]</sup>) so that select the best candidates in the decoding step is based on both conditional probability of the partial sequence outputted by the NMT model and the estimated long-term reward outputted by the value network.<br />
<br />
==Contributions==<br />
1) Developing a decoding scheme that considers long-term reward while generating words one by one for machine translation.<br />
<br />
2) Building another two modules for the value network, a semantic matching module and a context-coverage module. The semantic matching module estimates the similarity between the source and target sentences. The context-coverage module measures the coverage of context used in the encoder-decoder layer observing the fact that better translation is generated when using more context in the attention mechanism.<br />
<br />
<br />
Results of translation experiments demonstrate the effectiveness and robustness of the new decoding mechanism compared to several baseline algorithms.<br />
<br />
<br />
=Neural Machine Translation=<br />
<br />
NTM systems are implemented with a RNN based encoder-decoder framework, which directly models the probability $P(y|x)$ of a target sentence $y = {y_1,...,y_{T_y}}$ conditioned on the source sentence $x = {x_1,...,x_{T_x}}$, where $T_x$ and $T_y$ are the length of sentence x and y.<br />
<br />
The encoder of NMT reads the source sentence x word by word and generates a hidden representation for each word xi:<br />
$$<br />
h_i = f(h_{i-1},x_i)<br />
$$<br />
where function f is the recurrent unit such as LSTM unit or GRU.<br />
<br />
Then the decoder of NMT computes the conditional probability of each target word $y_t$ conditioned on its proceeding words y<t and source sentence:<br />
$$<br />
\begin{align*}<br />
c_t = q(r_{t-1};h_1,...,h_{T_x})\\<br />
r_t = g(r_{t-1},y_{t-1},c_t)\\<br />
P(y_t|y_{<t},x)\propto exp(y_t;r_t,c_t)<br />
\end{align*}<br />
$$<br />
where,<br />
<br />
$c_t$: weighted contextual information summarizing the source sentence x using some attention mechanism.<br />
<br />
$r_t$: decoder RNN hidden representation at step t, computed by an LSTM or GRU<br />
<br />
Notations:<br />
# $Θ$: All learned parameters<br />
# $π_{Θ}$: Translation model with parameter $Θ$<br />
# D: training dataset that contains source-target sentence pairs.<br />
The training process aims at seeking the optimal parameters $Θ^*$ using MLE to encode source sentence and decode it into the target sentence.<br />
<br />
<math>Θ^* = \displaystyle\arg\max_{Θ} \prod\limits_{(x,y)∈D}P(y|x;Θ) = \displaystyle\arg\max_{Θ}\prod\limits_{(x,y)∈D}\prod\limits^{T_y}_{t=1}P(y_t|y_{<t},x;Θ) </math><br />
<br />
=Value Network for NMT=<br />
'''Motivation''':<br />
Beam search has limited ability at predicting a high-quality sentence due to myopic bias. [3]<sup>[[#References|[3]]]</sup> introduced scheduled sampling approach, which takes the generated outputs from the model and the golden truth sentence in training to help the model learn from its own errors but cannot avoid the myopic bias of beam search during testing. [4]<sup>[[#References|[4]]]</sup> learns a predictor to predict the ranking score of a certain word at step t, and use this score to replace the conditional probability outputted by the NMT model for beam search during testing. Unfortunately, this work still looks only one step forward and cannot address the problem. Thus the authors motivate to estimate the expected performance of any sequence during decoding based on the concept of value function in reinforcement learning.<br />
<br />
==Value Network Structure==<br />
In conventional reinforcement learning, a value function describes how much cumulated reward could be collected from state s by following certain policy π. In NMT, consider $x$ and $y_{<t}$ as the state, $π_{Θ}$ as policy which can generate a word (action) given any state. The value function characterizes the expected translation performance (e.g. BLEU score) when using $π_{Θ}$ to translate $x$ with the previous t-1 words, $y_{<t}$. The value function is defined as:<br />
<br />
<math>v(x,y_{<t}) = \sum\limits_{y'∈Y: y'_{<t}=y_{<t}}BLEU(y^*(x),y')P(y'|x;Θ)</math><br />
<br />
where $y^*(x)$: ground truth translation, $Y$: the space of complete sentence.<br />
<br />
<br />
Two modules are developed to fully exploit the information in the encoder-decoder framework.<br />
[[File:VNFigure5.png|thumb|600px|center|Figure 5: Architecture of Value Network]]<br />
'''Semantic Matching (SM) Module''': at time step, use mean pooling over the decoder RNN hidden states and the over the context states as summarizations of the partial target sequence and the context in source language. Then use a feed-forward network to evaluate semantic information between the source sentence and the target sentence.<br />
i.e.,<br />
# $\bar{r_{t}} = \frac{1}{t}\sum\limits^{t}_{l=1}r_l$, $\bar{c_{t}} = \frac{1}{t}\sum\limits^{t}_{l=1}c_l$<br />
# $𝜇_{SM} = f_{SM}([\bar{r_{t}},\bar{c_{t}}])$<br />
<br />
'''Context-Coverage (CC) Module''': It is often observed that more context covered in the attention model leads to better translation. CC is built to measure coverage of information in the network. Similarly as the process in SM module, the process is defined as:<br />
# $\bar{h} = \frac{1}{T_x}\sum\limits^{T_x}_{l=1}h_l$<br />
# $𝜇_{CC} = f_{CC}([\bar{c_{t}},\bar{h}])$<br />
<br />
In the end, the authors concatenate $𝜇_{SM}$ and $𝜇_{CC}$ and then use another fully connected layer with sigmoid activation function to output a scalar as the value prediction. (Figure 5)<br />
<br />
==Training and Learning==<br />
The authors adopt the Monte-Carlo method to learn the value function. The training of the value network for NMT model $π_{Θ}$ is shown in Algorithm 1. <br />
[[File:VNFigure6.png|thumb|600px|center]]<br />
[[File:VNFigure7.png|thumb|600px|center]]<br />
[[File:VNFigure8.png|thumb|600px|center]]<br />
<br />
Notes:<br />
* 4-8 for training; 9 for learning<br />
* 9: Eqn.(7), For value function estimation, the authors minimize the pairwise ranking loss objective function to learn the mapping from $(x,y_p)$ → avg_bleu$(x,y_p)$ since they hope the value network to be useful in differentiating good and bad examples.<br />
* 11: $w$ is the learned weights of the value network<br />
<br />
==Decoding Process==<br />
<br />
The authors linearly combine the conditional probability $P(y|x)$, which is the output of the NMT model and the value network representing the future reward motivated by the success of AlphaGo<sup>[[#References|[5]]]</sup> They believe it would generate a better result by considering both past and future information. Mathematically,given a translation model $P(y|x)$, a value network $v(x,y)$ and a hyperparameter $𝛼∈(0,1)$, the '''score''' of partial sequence y for x is:<br />
<br />
$𝛼×\frac{1}{|y|}\,log\,P(y|x) + (1-𝛼)×log\,v(x,y)$<br />
<br />
where $|y|$ is the length of y. The details of the decoding process are presented in Algorithm 2. This neural network based decoding algorithm is called NMT-VNN for short.<br />
[[File:VNFigure9.png|thumb|600px|center]]<br />
<br />
=Experiments=<br />
==Settings==<br />
The author compare their proposed NMT-VNN with two baseline models, the classic NMT with beam search (NMT-BS)<sup>[[#References|[6]]]</sup> and the one referred as beam search optimization (NMT-BSO), which trains a predictor to evaluate the quality of any partial sequence, and then uses the predictor to select words instead of the probability. <br />
Note the difference between NMT-BSO and their NMT-VNN is NTM-BSO predicts the local improvement of BLEU for any single word, while NMT-VNN predicts the '''final''' BLEU score and use the predicted score to '''select words'''.<br />
<br />
The models are tested using three pairs of languages: English→French (En→Fr), English→German (En→De), and Chinese→English (Zh→En). They used the same bilingual corpora from WMT’14 as used in [6] , which contains 12M, 4.5M and 10M training data for each task.<sup>[[#References|[6]]]</sup><br />
* For En→Fr and En→De: validation set: newstest2012 and newstest2013; test set: newstest2014<br />
* For Zh→En: validation set: NIST 2004; test set: NIST 2006 and NIST 2008<br />
For all datasets in Chinese, they use a public tool for Chinese word segmentation. In all experiments, validation sets were only used for early-stopping and hyperparameter tuning.<br />
<br />
For NMT-VNN and NMT-BS, they first set experimental parameters to train an NMT model following <sup>[[#References|[6]]]</sup>. The vocabulary for each language is the most common 30K in the parallel corpora and the words not in the vocabulary (i.e., unknown words) were replaced with a special token “UNK". Each word was embedded into a vector space of 620 dimensions, and recurrent unit has dimension of 1000. Sentences with length ≤ 50 were kept in the training set. Batch size was set as 80 with 20 batches pre-fetched and sorted by sentence lengths.<br />
<br />
The NMT model was trained with asynchronized SGD on four K40m GPUs for about seven days. For NMT-BSO, they implemented the algorithm and the model was trained in the same environment.<br />
For the value network used in NMT-VNN, they set the same parameters for the encoder-decoder layers as the NMT model. Additionally, in the SM module and CC module, they set function $μ_{SM}$ and $μ_{CC}$ as single-layer feed forward networks with 1000 output nodes. In Algorithm 1, they set the hyperparameter K = 20 to estimate the value of any partial sequence. They used mini-batch training with batch size 80, and the value network model was trained with AdaDelta <sup>[[#References|[8]]]</sup> on one K40m GPU for about three days.<br />
<br />
<br />
During testing, the hyperparameter $𝛼$ for NMT-VNN was set by cross validation and they found the optimal $𝛼$ for En→Fr, En→De and Zh→En are 0.85, 0.9 and 0.8 respectively. They used the BLEU score <sup>[[#References|[9]]]</sup> as the evaluation metric, which is computed by the multi-bleu.perl<sup>[[#References|[10]]]</sup>. Beam search size was set to be 12 for all the algorithms following the common practice <sup>[[#References|[11]]]</sup>.<br />
<br />
==Results==<br />
[[File:VNFigure10.png|thumb|400px|right]]<br />
[[File:VNFigure11.png|thumb|400px|right|Figure 6: BLEU Score vs. Source Sentence Length]]<br />
Table 1 shows that the NMT-VNN algorithm outperforms the baseline algorithms on all tasks, especially for harder level translations (Zh→En). <br />
<br />
For En→Fr and En→De tasks, NMT-VNN outperforms NMT-BS by 1.03/1.3 points due to the additional use of value network, which suggests the additional knowledge provides useful information to help the NMT model. NMT-VNN outperforms NMT-BSO by about 0.31/0.33 points since NMT-BSO only uses a local BLEU predictor to estimate the partial BLEU score while NMT-VNN predicts the future performance, which shows the advantage of considering long-term rewards. For Zh→En task, NMT-VNN outperforms NMT-BS by 1.4/1.82 points on NIST 2006 and NIST 2008, and outperforms NMT-BSO by 1.01/0.72 points.<br />
<br />
The plots BLEU Score vs. Source Sentence Length (Figure 6) depict NMT-VNN algorithm outperforms the baseline algorithms for almost any sentence length.<br />
<br />
Furthermore, they also test the value network on a deep NMT model in which the encoder and decoder both have 4-layer LSTMs. The result (Table 1) shows 0.33 points improvement on the En→Fr task. These results demonstrate the effectiveness and robustness of the NMT-VNN algorithm.<br />
<br />
==Analysis on Value Network==<br />
[[File:VNFigure12.png|thumb|600px|center|Figure 7: (a)BLEU scores of En→De task w.r.t different beam size (b)BLEU scores of En→De<br />
task w.r.t different hyperparameter 𝛼]]<br />
First, the additional component in decoding will affect the translation process.<br />
<br />
Second, the accuracy of NMT sometimes drops dramatically as the beam size grows on certain tasks because the training of NMT favors short but inadequate translation candidates<sup>[[#References|[12]]]</sup>. This happens for the En→De translation, however, such shortage can be largely avoided by using value network. Figure 7(a) shows the accuracy (BLEU score) using different beam size for the NMT-BS and the NMT-VNN models. NMT-VNN is much more stable than the original NMT as its accuracy only differs a little for different beam sizes while NMT-BS drops more than 0.5 points when the beam size is large.<br />
<br />
Third, the performances of NMT-VNN using different hyperparameter 𝛼 during decoding for En→De task is shown in Figure 7(b) and it is stable for 𝛼 ranging from 0.7 to 0.95, and slightly drops for a smaller 𝛼. The authors declare the proposed algorithm is robust to the hyperparameter.<br />
<br />
=Conclusions and Future Work=<br />
This paper introduces a new decoding scheme that incorporates value networks for NMT. The new decoding scheme considers not only the local conditional probability of a candidate word but also its long-term reward for future decoding. Experiments on three translation tasks verify the effectiveness of the new scheme. For future works:<br />
*designing better structures for the value network.<br />
*extending it to other sequence-to-sequence learning tasks, such as image captioning and dialog systems.<br />
<br />
==Critiques==<br />
# It is a good idea to consider future reward by using the value network in addition to beam search which is based on past information. Intuitively, the NTM-VNN should improve the NMT by using both past and future information.<br />
# The paper does not give much information or give any quantitative evaluations about the two modules they built, for example, how much the two modules contribute to the accuracy of their NMT-VNN model?<br />
# In algorithm 1 step 5, it is reasonable to generate partial target sentence $y_{p,1}$,$y_{p,2}$ using $π_{Θ}$ with random early stop, but if the stop times for the two sentences are too far from each other, there could be a problem for the beam search...since beam size is kind of related to the number of words of the sentence left to generate.<br />
# In the experiment results section, they only test the value network on a deep NMT model on the En→Fr task and show the improvement, is that true for all other translations?<br />
#In the decoder section of the network, the authors could have experimented with fast-forward linear connections while stacking LSTMs. This technique has proven to obtain some of the best empirical accuracies in machine translation ( En->Fr BLEU= 37.7) [13].<br />
<br />
=References=<br />
<br />
[1] https://github.com/tensorflow/nmt<br />
<br />
[2] https://en.wikipedia.org/wiki/BLEU<br />
<br />
[3] S. Bengio, O. Vinyals, N. Jaitly, and N. Shazeer. Scheduled sampling for sequence prediction with recurrent neural networks. In Advances in Neural Information Processing Systems, pages 1171–1179, 2015.<br />
<br />
[4] S. Wiseman and A. M. Rush. Sequence-to-sequence learning as beam-search optimization. In EMNLP, 2016.<br />
<br />
[5] D. Silver, A. Huang, C. J. Maddison, A. Guez, L. Sifre, G. Van Den Driessche, J. Schrittwieser, I. Antonoglou, V. Panneershelvam, M. Lanctot, et al. Mastering the game of go with deep neural networks and tree search. Nature, 529(7587):484–489, 2016.<br />
<br />
[6] D. Bahdanau, K. Cho, and Y. Bengio. Neural machine translation by jointly learning to align and translate. ICLR, 2015.<br />
<br />
[7] [https://baike.baidu.com/item/%E4%B8%AD%E6%96%87%E5%88%87%E8%AF%8D Chinese word segmentation]<br />
<br />
[8] M. D. Zeiler. Adadelta: an adaptive learning rate method. arXiv preprint arXiv:1212.5701, 2012.<br />
<br />
[9] K. Papineni, S. Roukos, T. Ward, and W.-J. Zhu. Bleu: a method for automatic evaluation of machine translation. In Proceedings of the 40th annual meeting on association for computational linguistics, pages 311–318. Association for Computational Linguistics, 2002.<br />
<br />
[10] https://github.com/moses-smt/mosesdecoder/blob/master/scripts/generic/multi-bleu.perl.<br />
<br />
[11] I. Sutskever, O. Vinyals, and Q. V. Le. Sequence to sequence learning with neural networks. In Advances in neural information processing systems, pages 3104–3112, 2014.<br />
<br />
[12] Z. Tu, Y. Liu, L. Shang, X. Liu, and H. Li. Neural machine translation with reconstruction. In AAAI, pages 3097–3103, 2017.<br />
<br />
[13] Jie Zhou, Ying Cao, Xuguang Wang, Peng Li, Wei Xu ; "Deep Recurrent Models with Fast-Forward Connections for Neural Machine Translation". arXiv:1606.04199 [cs.CL]</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/Conditional_Image_Generation_with_PixelCNN_Decoders&diff=31295STAT946F17/Conditional Image Generation with PixelCNN Decoders2017-11-23T18:13:54Z<p>Nkadiyar: </p>
<hr />
<div>=Introduction=<br />
This works is based of the widely used PixelCNN and PixelRNN, introduced by Oord et al. in [[#Reference|[1]]]. From the previous work, the authors observed that PixelRNN performed better than PixelCNN, however, PixelCNN was faster to compute as you can parallelize the training process. In this work, Oord et al. [[#Reference|[2]]] introduced a Gated PixelCNN, which is a convolutional variant of the PixelRNN model, based on PixelCNN. In particular, the Gated PixelCNN uses explicit probability densities to generate new images using autoregressive connections to model images through pixel-by-pixel computation by decomposing the joint image distribution as a product of conditionals [[#Reference|[6]]]. The Gated PixelCNN is an improvement over the PixelCNN by removing the "blindspot" problem, and to yield a better performance, the authors replaced the ReLU units with sigmoid and tanh activation function. The proposed Gated PixelCNN combines the strength of both PixelRNN and PixelCNN - that is by matching the log-likelihood of PixelRNN on both CIFAR and ImageNet along with the quicker computational time presented by the PixelCNN [[#Reference|[9]]]. Moreover, the authors also introduced a conditional Gated PixelCNN variant (called Conditional PixelCNN) which has the ability to generate images based on class labels, tags, as well as latent embedding to create new image density models. These embeddings capture high level information of an image to generate a large variety of images with similar features; for instance, the authors can generate different poses of a person based on a single image by conditioning on a one-hot encoding of the class. This approach provided insight into the invariances of the embeddings which enabled the authors to generate different poses of the same person based on a single image. Finally, the authors also presented a PixelCNN Auto-encoder variant which essentially replaces the deconvolutional decoder with the PixelCNN.<br />
<br />
=Gated PixelCNN=<br />
Pixel-by-pixel is a simple generative method wherein given an image of dimension of dimension $x_{n^2}$, we iterate, employ feedback and capture pixel densities from every pixel to predict our "unknown" pixel density $x_i$. To do this, the traditional PixelCNNs and PixelRNNs adopted the joint distribution p(x), wherein the pixels of a given image is the product of the conditional distributions. Hence, the authors employ autoregressive models which means they just use plain chain rule for joint distribution, depicted in Equation 1. So the very first pixel is independent, second depend on first, third depends on first and second and so on. Basically you just model your image as sequence of points where each pixel depends linearly on previous ones. Equation 1 depicts the joint distribution where x_i is a single pixel:<br />
<br />
$$p(x) = \prod\limits_{i=1}^{n^2} p(x_i | x_1, ..., x_{i-1})$$<br />
<br />
where $p(x)$ is the generated image, $n^2$ is the number of pixels, and $p(x_i | x_1, ..., x_{i-1})$ is the probability of the $i$th pixel which depends on the values of all previous pixels. It is important to note that $p(x_0, x_1, ..., x_{n^2})$ is the joint probability based on the chain rule - which is a product of all conditional distributions $p(x_0) \times p(x_1|x_0) \times p(x_2|x_1, x_0)$ and so on. Figure 1 provides a pictorial understanding of the joint distribution which displays that the pixels are computed pixel-by-pixel for every row, and the forthcoming pixel depends on the pixels values above and to the left of the pixel in concern. <br />
<br />
[[File:xi_img.png|500px|center|thumb|Figure 1: Computing pixel-by-pixel based on joint distribution.]]<br />
<br />
Hence, for every pixel, we use the softmax layer towards the end of the PixelCNN to predict the pixel intensity value (i.e. the highest probable index from 0 to 255). Figure 2 [[#Reference|[7]]] illustrates how to predict (generate) a single pixel value.<br />
<br />
[[File:single_pixel.png|500px|center|thumb|Figure 2: Predicting a single pixel value based on softmax layer.]]<br />
<br />
So, the PixelCNN is supposedly to maps a neighborhood of pixels to prediction for the next pixel. That is, to generate pixel $x_i$ the model can only condition on the previously generated pixels $x_1 , ..., x_{i−1}$; so every conditional distribution is modelled by a convolutional neural network. For instance, given a $5\times5$ image (let's represent each pixel as an alphabet and zero-padded), and we have a filter of dimension $3\times3$ that slides over the image which multiplies each element and sums them together to produce a single response. However, we cannot use this filter because pixel $a$ should not know the pixel intensities for $b,f,g$ (future pixel values). To counter this issue, the authors use a mask on top of the filter to only choose prior pixels and zeroing the future pixels to negate them from calculation - depicted in Figure 3 [[#Reference|[7]]]. Hence, to make sure the CNN can only use information about pixels above and to the left of the current pixel, the filters of the convolution are masked - that means the model cannot read pixels below (or strictly to the right) of the current pixel to make its predictions [[#Reference|[3]]]. <br />
<br />
[[File:masking1.png|200px|center|thumb|Figure 3: Masked convolution for a $3\times3$ filter.]]<br />
[[File:masking2.png|500px|center|thumb|Figure 4: Masked convolution for each convolution layer.]]<br />
<br />
Hence, for each pixel there are three colour channels (R, G, B) which are modeled successively, with B conditioned on (R, G), and G conditioned on R [[#Reference|[8]]]. This is achieved by splitting the feature maps at every layer of the network into three and adjusting the centre values of the mask tensors, as depicted in Figure 5 [[#Reference|[8]]]. The 256 possible values for each colour channel are then modelled using a softmax.<br />
<br />
[[File:rgb_filter.png|300px|right|thumb|Figure 5: RGB Masking.]]<br />
<br />
Now, from Figure 6, notice that as the filter with the mask slides across the image, pixel $f$ does not take pixels $c, d, e$ into consideration (breaking the conditional dependency) - this is where we encounter the "blind spot" problem. <br />
<br />
[[File:blindspot.gif|500px|center|thumb|Figure 6: The blindspot problem.]]<br />
<br />
It is evident that the progressive growth of the receptive field of the masked kernel over the image disregards a significant portion of the image. For instance, when using a 3x3 filter, roughly quarter of the receptive field is covered by the "blind spot", meaning that the pixel contents are ignored in that region. In order to address the blind spot, the authors use two filters (horizontal and vertical stacks) in conjunction to allow for capturing the whole receptive ﬁeld, depicted in Figure{vh_stack}. In particular, the horizontal stack conditions the current row, and the vertical stack conditions all the rows above the current pixel. It is observed that the vertical stack, which does not have any masking, allows the receptive ﬁeld to grow in a rectangular fashion without any blind spot. Thereafter, the outputs of both the stacks, per-layer, is combined to form the output. Hence, every layer in the horizontal stack takes an input which is the output of the previous layer as well as that of the vertical stack. By splitting the convolution into two different operations enables the model to access all pixels prior to the pixel of interest. <br />
<br />
[[File:vh_stack.png|500px|center|thumb|Figure 7: Vertical and Horizontal stacks.]]<br />
<br />
=== Horizontal Stack ===<br />
For the horizontal stack (in purple for Figure 7), the convolution operation conditions only on the current row, so it has access to left pixels. In essence, we take a $1 \times n//2+1$ convolution with shift (pad and crop) rather than $1\times n$ masked convolution. So, we perform convolution on the row with a kernel of width 2 pixels (instead of 3) from which the output is padded and cropped such that the image shape stays the same. Hence, the image convolves with kernel width of 2 and without masks.<br />
<br />
[[File:h_mask.png|500px|center|thumb|Figure 8: Horizontal stack.]]<br />
<br />
Figure 8 [[#Reference|[3]]] shows that the last pixel from output (just before ‘Crop here’ line) does not hold information from last input sample (which is the dashed line).<br />
<br />
=== Vertical Stack ===<br />
Vertical stack (blue) has access to all top pixels. The vertical stack is of kernel size $n//2 + 1 \times n$ with the input image being padded with another row in the top and bottom. Thereafter, we perform the convolution operation, and crop the image to force the predicted pixel to be dependent on the upper pixels only (i.e. to preserve the spatial dimensions)[[#Reference|[3]]]. Since the vertical filter does not contain any "future" pixel values, only upper pixel values, no masking is incorporated as no target pixel is touched. However, the computed pixel from the vertical stack yields information from top pixels and sends that info to horizontal stack (which supposedly eliminates the "blindspot problem").<br />
<br />
[[File:vertical_mask.gif|500px|center|thumb|Figure 9: Vertical stack.]]<br />
<br />
From Figure 9 it is evident that the image is padded (left) with kernel height zeros, then convolution operation is performed from which we crop the output so that rows are shifted by one with respect to input image. Hence, it is noticeable that the first row of output does not depend on first (real, non-padded) input row. Also, the second row of output only depends on the first input row - which is the desired behaviour.<br />
<br />
=== Gated block ===<br />
The PixelRNNs are observed to perform better than the traditional PixelCNN for generating new images. This is because the spatial LSTM layers in the PixelRNN allows for every layer in the network to access the entire neighborhood of previous pixels. The PixelCNN, however, only takes into consideration the neighborhood region and the depth of the convolution layers to make its predictions [[#Reference|[4]]]. Another advantage for the PixelRNN is that this network contains multiplicative units (in the form of the LSTM gates), which may help it to model more complex interactions [[#Reference|[3]]]. To address the benefits of PixelRNN and append it onto the newly proposed Gated PixelCNN, the authors replaced the rectified linear units between the masked convolutions with the following custom-made gated activation function, depicted in Equation 2:<br />
<br />
$$y = tanh(W_{k,f} \ast x) \odot \sigma(W_{k,g} \ast x)$$<br />
<br />
where $\sigma$ is the sigmoid non-linearity, $k$ is the number of the layer, $f, g$ are the different feature maps, $\odot$ is the element-wise product and $\ast$ is the convolution with the input. This function is the key ingredient that cultivates the Gated PixelCNN model. <br />
<br />
Figure 10 provides a pictorial illustration of a single layer in the Gated PixelCNN architecture; wherein the vertical stack contributes to the horizontal stack with the $1\times1$ convolution - going the other way would break the conditional distribution. In other words, the horizontal and vertical stacks are sort of independent, wherein vertical stack should not access any information horizontal stack has - otherwise it will have access to pixels it shouldn’t see. However, vertical stack can be connected to vertical as it predicts pixel following those in vertical stack. In particular, the convolution operations are shown in green (which are masked), element-wise multiplications and additions are shown in red. The convolutions with $W_f$ and $W_g$ are not combined into a single operation (which is essentially the masked convolution) to increase parallelization shown in blue. The parallelization now splits the $2p$ features maps into two groups of $p$. Finally, the authors also use the residual connection in the horizontal stack. Moreover, the $(n \times 1)$ and $(n \times n)$ are the masked convolutions which can also be implemented as $([n/2] \times 1)$ and $([n/2] \times n)$ which are convolutions followed by a shift in pixels by padding and cropping to get the original dimension of the image.<br />
<br />
[[File:gated_block.png|500px|center|thumb|Figure 10: Gated block.]]<br />
<br />
In essence, PixelCNN typically consists of a stack of masked convolutional layers that takes an $N \times N \times 3$ image as input and produces $N \times N \times 3 \times 256$ (probability of pixel intensity) predictions as output. During sampling the predictions are sequential: every time a pixel is predicted, it is fed back into the network to predict the next pixel. This sequentiality is essential to generating high quality images, as it allows every pixel to depend in a highly non-linear and multimodal way on the previous pixels. <br />
<br />
Another important mention is that the residual connections are only for horizontal stacks. On the other side skip connections allow as to incorporate features from all layers at the very end of out network. Most important stuff to mention here is that skip and residual connection use different weights after gated block.<br />
<br />
=Conditional PixelCNN=<br />
Conditioning is a smart word for saying that we’re feeding the network some high-level information - for instance, providing an image to the network with the associated classes in MNIST/CIFAR datasets. During training you feed image as well as class to your network to make sure network would learn to incorporate that information as well. During inference you can specify what class your output image should belong to. You can pass any information you want with conditioning, we’ll start with just classes.<br />
<br />
For a conditional PixelCNN, we represent a provided high-level image description as a latent vector $h$, wherein the purpose of the latent vector is to model the conditional distribution $p(x|h)$ such that we get a probability as to if the images suites this description. The conditional PixelCNN models based on the following distribution:<br />
<br />
$$p(x|h) = \prod\limits_{i=1}^{n^2} p(x_i | x_1, ..., x_{i-1}, h)$$<br />
<br />
Hence, now the conditional distribution is dependent on the latent vector h, which is now appended onto the activations prior to the non-linearities; hence the activation function after adding the latent vector becomes:<br />
<br />
$$y = tanh(W_{k,f} \ast x + V_{k,f}^T h) \odot \sigma(W_{k,g} \ast x + V_{k,g}^T h)$$<br />
<br />
Note $h$ multiplied by matrix inside tanh and sigmoid functions, $V$ matrix has the shape [number of classes, number of filters], $k$ is the layer number, and the classes were passed as a one-hot vector $h$ during training and inference.<br />
<br />
Note that if the latent vector h is a one-hot encoding vector that provides the class labels, which is equivalent to the adding a class dependent bias at every layer. So, this means that the conditioning is independent from the location of the pixel - this is only if the latent vector holds information about “what should the image contain” rather than the location of contents in the image. For instance, we could specify that a certain animal or object should appear in different positions, poses and backgrounds.<br />
<br />
In addition, the authors also developed a variant that makes the conditional distribution dependent on the location (an application when the location of an object is important). This is achieved by mapping the latent vector $h$ to a spatial representation $s=m(h)$ (which contains the same dimension of the image but may have an arbitrary number of feature maps) with a deconvolutional neural network $m()$; this provides a location dependent bias as follows:<br />
<br />
$$y = tanh(W_{k,f} \ast x + V_{k,f} \ast s) \odot \sigma(W_{k,g} \ast x + V_{k,g} \ast s)$$<br />
<br />
where $V_{k,g}\ast s$ is an unmasked $1\times1$ convolution.<br />
<br />
=== PixelCNN Auto-Encoders ===<br />
Since conditional PixelCNNs can model images based on the distribution $p(x|h)$, it is possible to apply this analogy into image decoders used in Autoencoders. Introduced by Hinton et. al in [[#Reference|[5]]], autoencoder is a dimensionality reduction neural network which is composed of two parts: an encoder which maps the input image into low-dimensional representation (i.e. the latent vector $h$) , and a decoder that decompresses the latent vector to reconstruct the original image. <br />
<br />
In order to apply the conditional PixelCNN onto the autoencoder, the deconvolutional decoders are replaced with the conditional PixelCNN - the re-architectured network of which is used for training a data set. The authors observe that the encoder can better extract representations of the provided input data - this is because much of the low-level pixel statistics is now handled by the PixelCNN; hence, the encoder omits low-level pixel statistics and focuses on more high-level abstract information.<br />
<br />
=Experiments=<br />
<br />
===Unconditional Modelling with Gated PixelCNN===<br />
For the first set of experiments, the authors evaluate the Gated PixelCNN unconditioned model on the CIFAR-10 dataset is adopted. A comparison of the validation score between the Gated PixelCNN, PixelCNN, and PixelRNN is computed, wherein the lower score means that the optimized model generalizes better. Using the negative log-likelihood criterion (NLL), the Gated PixelCNN obtains an NLL Test (Train) score of 3.03 (2.90) which outperforms the PixelCNN by 0.11 bits/dim, which obtains 3.14 (3.08). Although the performance is a bit better, visually the quality of the samples that were produced is much better for the Gated PixelCNN when compared to PixelCNN. It is important to note that the Gated PixelCNN came close to the performance of PixelRNN, which achieves a score of 3.00 (2.93). Table 1 provides the test performance of benchmark models on CIFAR-10 in bits/dim (where lower is better), and the corresponding training performance is in brackets.<br />
<br />
[[File:ucond_cifar.png|500px|center|thumb|Table 1: Evaluation on CIFAR-10 dataset for an unconditioned GatedPixelCNN model.]]<br />
<br />
Another experiment on the ImageNet data is performed for image sizes $32 \times 32$ and $64 \times 64$. In particular, for a $32 \times 32$ image, the Gated PixelCNN obtains a NLL Test (Train) of 3.83 (3.77) which outperforms PixelRNN which achieves 3.86 (3.83); from which the authors observe that larger models do have better performance, however, the simpler PixelCNN does have the ability to scale better. For a $64 \times 64$ image, the Gated PixelCNN obtains 3.57 (3.48) which, yet again, outperforms PixelRNN which achieves 3.63 (3.57). The authors do mention that the Gated PixelCNN performs similarly to the PixelRNN (with row LSTM); however, Gated PixelCNN is observed to train twice as quickly at 60 hours when using 32 GPUs. The Gated PixelCNN has 20 layers (Figure 2), each of which has 384 hidden units and a filter size of 5x5. For training, a total of 200K synchronous updates were made over 32 GPUs which were computed in TensorFlow using a total batch size of 128. Table 2 illustrates the performance of benchmark models on ImageNet dataset in bits/dim (where lower is better), and the training performance in brackets.<br />
<br />
[[File:ucond_imagenet.png|500px|center|thumb|Table 1: Evaluation on ImageNet dataset for an unconditioned GatedPixelCNN model.]]<br />
<br />
<br />
===Conditioning on ImageNet Classes===<br />
For the the second set of experiments, the authors evaluated the Gated PixelCNN model by conditioning the classes of the ImageNet images. Using the one-hot encoding $(h_i)$, for which the $i^th$ class the distribution becomes $p(x|h_i)$, the model receives roughly log(1000) $\approx$ 0.003 bits/pixel for a $32 \times 32$ image. Although the log-likelihood did not show a significant improvement, visually the quality of the images were generated much better when compared to the original PixelCNN. <br />
<br />
Figure 11 shows some samples from 8 different classes of ImageNet images from a single class-conditioned model. It is evident that the Gated PixelCNN can better distinguish between objects, animals and backgrounds. The authors observe that the model can generalize and generate new renderings from the animal and object class, when the trained model is provided with approximately 1000 images.<br />
<br />
[[File:cond_imagenet.png|500px|center|thumb|Figure 11: Class-Conditional samples from the Conditional PixelCNN on the ImageNet dataset.]]<br />
<br />
<br />
===Conditioning on Portrait Embeddings===<br />
For the third set of experiments, the authors used the top layer of the CNN trained on a large database of portraits that were automatically cropped from Flickr images using face detector. This pre-trained network was trained using triplet loss function which ensured a similar the latent embeddings for particular face across the entire dataset. <br />
<br />
In essence, the authors took the latent vector from this supervised pre-trained network which now has the architecture (image=$x$, embedding=$h$) tuples and trained the<br />
Conditional PixelCNN with the latent embeddings to model the distribution $p(x|h)$. Hence, if the network is provided with a face that is not in the training set, the model now has the capability to compute the latent embeddings $h=f(x)$ such that the output will generate new portraits of the same person. Figure 12 provides a pictorial example of the aforementioned manipulated network where it is evident that the generative model can produce a variety of images, independent from pose and lighting conditions, by extracting the latent embeddings from the pre-trained network. <br />
<br />
[[File:cond_portrait.png|500px|center|thumb|Figure 12: Input image is to the lest, whereas the portraits to the right are generated from high-level latent representation.]]<br />
<br />
<br />
===PixelCNN Auto Encoder===<br />
For the final set of experiment, the authors venture the possibility to train the a Gated PixelCNN by adopting the Autoencoder architecture. The authors start by training a PixelCNN auto-encoder using $32 \times 32$ ImageNet patches and compared its results to a convolutional autoencoder, optimized using mean-square error. It is important to note that both the models use a 10 or 100 dimensional bottleneck. <br />
<br />
Figure 13 provides a reconstruction using both the models. It is evident that the latent embedding produced when using PixelCNN autoencoder is much different when compared to convolutional autoencoder. For instance, in the last row, the PixelCNN autoencoder is able to generate similar looking indoor scenes with people without directly trying to "reconstruct" the input, as done by the convolutional autoencoder.<br />
<br />
[[File:pixelauto.png|500px|center|thumb|Figure 13: From left to right: original input image, reconstruction by an autoencoder trained with MSE, conditional samples from a PixelCNN as the deconvolution to the autoencoder. It is important to note that both these autoencoders were trained end-to-end with 10 and 100-dimensional bottleneck values.]]<br />
<br />
<br />
=Conclusion=<br />
This work introduced the Gated PixelCNN which is an improvement over the original PixelCNN. In addition to the Gated PixelCNN being more computationally efficient, it now has the ability to match, and in some cases, outperform PixelRNN. In order to deal with the "blind spots" in the receptive fields presented in the PixelCNN, the newly proposed Gated PixelCNN use two CNN stacks (horizontal and vertical filters) to deal with this problem. Moreover, the authors now use a custom-made tank and sigmoid function over the ReLU activation functions because these multiplicative units helps to model more complex interactions. The proposed network obtains a similar performance to PixelRNN on CIFAR-10, however, it is now state-of-the-art on the ImageNet $32 \times 32$ and $64 \times 64$ datasets. <br />
<br />
In addition, the conditional PixelCNN is also explored on natural images using three different settings. When using class-conditional generation, the network showed that a single model is able to generate diverse and realistic looking images corresponding to different classes. When looking at generating human portraits, the model does have the ability to generate new images from the same person in different poses and lightning conditions given a single image. Finally, the authors also showed that the PixelCNN can be used as image decoder in an autoencoder. Although the log-likelihood is quite similar when comparing it to literature, the samples generated from the PixelCNN autoencoder model does provide a high visual quality images showing natural variations of objects and lighting conditions.<br />
<br />
<br />
=Summary=<br />
$\bullet$ Improved PixelCNN called Gated PixelCNN<br />
# Similar performance as PixelRNN, and quick to compute like PixelCNN (since it is easier to parallelize)<br />
# Fixed the "blind spot" problem by introducing 2 stacks (horizontal and vertical)<br />
# Gated activation units which now use sigmoid and tanh instead of ReLU units<br />
<br />
$\bullet$ Conditioned Image Generation<br />
# One-shot conditioned on class-label<br />
# Conditioned on portrait embedding<br />
# PixelCNN AutoEncoders<br />
<br />
=Critique=<br />
# The paper is not descriptive, and does not explain well on how the horizontal and vertical stacks solve the "blindspot" problem. In addition, the authors just mention the "gated block" and how they designed it, but they do not explain the intuition and how this approach is an improvement over the PixelCNN <br />
# The authors do not provide a good pictorial representation on any of the aforementioned novelties<br />
# The PixelCNN AutoEncoder is not descriptive enough! <br />
<br />
=Reference=<br />
# Aaron van den Oord et al., "Pixel Recurrent Neural Network", ICML 2016<br />
# Aaron van den Oord et al., "Conditional Image Generation with PixelCNN Decoders", NIPS 2016<br />
# S. Turukin, "Gated PixelCNN", Sergeiturukin.com, 2017. [Online]. Available: http://sergeiturukin.com/2017/02/24/gated-pixelcnn.html. [Accessed: 15- Nov- 2017].<br />
# S. Reed, A. van den Oord, N. Kalchbrenner, V. Bapst, M. Botvinick and N. Freitas, "Generating interpretable images with controllable structure", 2016.<br />
# G. Hinton, "Reducing the Dimensionality of Data with Neural Networks", Science, vol. 313, no. 5786, pp. 504-507, 2006.<br />
# "Conditional Image Generation with PixelCNN Decoders", Slideshare.net, 2017. [Online]. Available: https://www.slideshare.net/suga93/conditional-image-generation-with-pixelcnn-decoders. [Accessed: 18- Nov- 2017].<br />
# "Gated PixelCNN", Kawahara.ca, 2017. [Online]. Available: http://kawahara.ca/conditional-image-generation-with-pixelcnn-decoders-slides/gated-pixelcnn/. [Accessed: 17- Nov- 2017].<br />
# K. Dhandhania, "PixelCNN + PixelRNN + PixelCNN 2.0 — Commonlounge", Commonlounge.com, 2017. [Online]. Available: https://www.commonlounge.com/discussion/99e291af08e2427b9d961d41bb12c83b. [Accessed: 15- Nov- 2017].<br />
# S. Turukin, "PixelCNN", Sergeiturukin.com, 2017. [Online]. Available: http://sergeiturukin.com/2017/02/22/pixelcnn.html. [Accessed: 17- Nov- 2017].</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Hierarchical_Question-Image_Co-Attention_for_Visual_Question_Answering&diff=31293Hierarchical Question-Image Co-Attention for Visual Question Answering2017-11-23T18:11:05Z<p>Nkadiyar: </p>
<hr />
<div>__TOC__<br />
== Paper Summary ==<br />
{| class="wikitable"<br />
|-<br />
|'''Conference'''<br />
| <br />
* NIPS 2016<br />
* Presented as spotlight oral: [https://www.youtube.com/watch?v=m6t9IFdk0ms Youtube link]<br />
* 85 citations so far<br />
|-<br />
| '''Authors'''<br />
|Jiasen Lu, Jianwei Yang, Dhruv Batra, '''Devi Parikh'''<br />
|-<br />
|'''Abstract'''<br />
|''A number of recent works have proposed attention models for Visual Question Answering (VQA) that generate spatial maps highlighting image regions relevant to answering the question. In this paper, we argue that in addition to modeling "where to look" or visual attention, it is equally important to model "what words to listen to" or question attention. We present a novel co-attention model for VQA that jointly reasons about image and question attention. In addition, our model reasons about the question (and consequently the image via the co-attention mechanism) in a hierarchical fashion via a novel 1-dimensional convolution neural networks (CNN). Our model improves the state-of-the-art on the VQA dataset from 60.3% to 60.5%, and from 61.6% to 63.3% on the COCO-QA dataset. By using ResNet, the performance is further improved to 62.1% for VQA and 65.4% for COCO-QA.''<br />
|}<br />
= Introduction =<br />
'''Visual Question Answering (VQA)''' is a recent problem in computer vision and<br />
natural language processing that has garnered a large amount of interest from<br />
the deep learning, computer vision, and natural language processing communities.<br />
In VQA, an algorithm needs to answer text-based questions about images in<br />
natural language as illustrated in Figure 1.<br />
<br />
[[File:vqa-overview.png|thumb|600px|center|Figure 1: Illustration of VQA system whereby machine learning algorithm answers a visual question asked by an user for a given image (ref: http://www.visualqa.org/static/img/challenge.png)]]<br />
<br />
Recently, ''visual-attention'' based models have gained traction for VQA tasks, where the<br />
attention mechanism typically produces a spatial map highlighting image regions<br />
relevant for answering the visual question about the image. However, to correctly answer the <br />
question, machine not only needs to understand or "attend"<br />
regions in the image but also the parts of question as well. In this paper, authors have proposed a novel ''co-attention''<br />
technique to combine "where to look" or visual-attention along with "what words<br />
to listen to" or question-attention VQA allowing their model to jointly reasons about image and question thus improving <br />
upon existing state of the art results.<br />
<br />
== "Attention" Models ==<br />
You may skip this section if you already know about "attention" in<br />
context of deep learning. Since this paper talks about "attention" almost<br />
everywhere, I decided to put this section to give very informal and brief<br />
introduction to the concept of the "attention" mechanism specially visual "attention", <br />
however, it can be expanded to any other type of "attention".<br />
<br />
Visual attention in CNN is inspired by the biological visual system. As humans,<br />
we have ability to focus our cognitive processing onto a subset of the<br />
environment that is more relevant for the given situation. Imagine, you witness<br />
a bank robbery where robbers are trying to escape on a car, as a good citizen,<br />
you will immediately focus your attention on number plate and other physical<br />
features of the car and robbers in order to give your testimony later, however, you may not remember things which otherwise interests you more. <br />
Such selective visual attention for a given context (robbery in above example) can also be implemented in<br />
traditional CNNs as well. This allows CNNs to be more robust and superior for certain tasks and it even helps <br />
algorithm designer to visualize what spacial features (regions within image) were more important than others. Attention guided<br />
deep learning is particularly very helpful for image caption and VQA tasks.<br />
<br />
== Role of Visual Attention in VQA ==<br />
This section is not a part of the actual paper that is been summarized, however, it gives an overview<br />
of how visual attention can be incorporated in training of a network for VQA tasks, eventually, helping <br />
readers to absorb and understand actual proposed ideas from the paper more effortlessly.<br />
<br />
Generally for implementing attention, network tries to learn the conditional <br />
distribution $P_{i \in [1,n]}(Li|c)$ representing individual importance for all the features <br />
extracted from each of the dsicrete $n$ locations within the image <br />
conditioned on some context vector $c$. In order words, given $n$ features <br />
$L_i = [L_1, ..., L_n]$ from $n$ different spacial regions within the image (top-left, top-middle, top-right, and so on), <br />
then "attention" module learns a parameteric function $F(c;\theta)$ that outputs an importance mapping <br />
of each of these individual feature for a given context vector $c$ or a discrete probability distribution <br />
of size $n$, can be achived by $softmax(n)$. <br />
<br />
In order to incorporate the visual attention in VQA task, one can define context vector $c$ <br />
as a representation of the visual question asked by an user (using RNN perhaps LSTM). The context $c$ can then be used to generate an <br />
attention map for corresponding image locations (as shown in Figure 2) further improving the accuracy on final end-to-end training. <br />
Most work that exists in literature regarding use of visual-attention in VQA tasks are generally further <br />
specialization of the similar ideas.<br />
<br />
[[File:attention-vqa-general.png|thumb|700px|center|Figure 2: Different attention maps generated based on the given visual question. Regions with most "attention" or importance is whitened, machine learning model has learned to steer its attention based on the given question.]]<br />
<br />
== Motivation and Main Contributions ==<br />
So far, all attention models for VQA in literature have focused on the problem of identifying "where<br />
to look" or visual attention. In this paper, authors argue that the problem of identifying "which words to<br />
listen to" or '''question attention''' is equally important. Consider the questions "how many horses are<br />
in this image?" and "how many horses can you see in this image?". They have the same meaning,<br />
essentially captured by the first three words. A machine that attends to the first three words would<br />
arguably be more robust to linguistic variations irrelevant to the meaning and answer of the question.<br />
Motivated by this observation, in addition to reasoning about visual attention, paper has addressed the<br />
problem of question attention. Basically, main contributions of the paper are as follows.<br />
<br />
* A novel co-attention mechanism for VQA that jointly performs question-guided visual attention and image-guided question attention.<br />
* A hierarchical architecture to represent the question, and consequently construct image-question co-attention maps at 3 different levels: word level, phrase level and question level. <br />
* A novel convolution-pooling strategy at phase-level to adaptively select the phrase sizes whose representations are passed to the question level representation.<br />
* Results on VQA and COCO-QA and ablation studies to quantify the roles of different components in the model<br />
<br />
= Method =<br />
This section is broken down into four parts: '''(i)''' notations used within the paper and also throughout this summary, '''(ii)''' hierarchical representation for a visual question, '''(iii)''' the proposed co-attention mechanism and<br />
'''(iv)''' predicting answers.<br />
<br />
== Notations ==<br />
{| class="wikitable"<br />
|-<br />
|'''Notation'''<br />
|'''Explaination'''<br />
|-<br />
|$Q = \{q_1,...q_T\}$<br />
|One-hot encoding of a visual question with $T$ words. Paper uses three different representation of visual question, one for each level of hierarchy, they are as follows: <br />
# $Q^w = \{q^w_1,...q^w_T\}$: Word level representation of visual question<br />
# $Q^p = \{q^p_1,...q^p_T\}$: Phrase level representation of visual question<br />
# $Q^s = \{q^s_1,...q^s_T\}$: Question level representation of visual question<br />
$Q^{w,p,s}$ has exactly $T$ number of embeddings in it (sequential data with temporal dimension), regardless of its position in the hierarchy i.e. word, phrase or question. <br />
|-<br />
|$V = {v_1,..,v_N}$<br />
|$V$ represented various vectors from $N$ different locations within the given image. Therefore, $v_n$ is feature vector from the image at location $n$. $V$ collectively covers the entire spatial reachings of the image. One can extract these location sensitive features from convolution layer of CNN.<br />
|-<br />
|$\hat{v}^r$ and $\hat{q}^r$<br />
|The co-attention features of image and question at each level in the hierarchy where $r \in \{w,p,s\}$. Basically, its a sum of $Q$ or $V$ after the dot product with attention $a^q$ or $a^v$ at each level of hierarchy. <br />
For example, at "word" level, $a^q_w$ and $a^v_w$ is a probability distribution representing importance of each words in visual question and each location within image respectively, whereas $\hat{q}^w$ and $\hat{v}^w$ are final features vectors for the given question and image with attention maps ($a^q_w$ and $a^v_w$ applied) at the "word" level, and similarly for "phrase" and "question" level as well.<br />
|}<br />
'''Note:''' Throughout the paper, $W$ represents the learnable weights and biases are not used within the equations for simplicity (reader must assume it to exist).<br />
<br />
== Question Hierarchy ==<br />
There are three levels of granularities for their hierarchical representation of a visual question: '''(i)''' word, '''(ii)''' phrase and '''(iii)''' question level. It is important to note, each level depends on the previous one, so, phrase level representations are extracted from word level and question level representations come from phrase level as depicted in Figure 4.<br />
<br />
[[File:hierarchy2.png|thumb|Figure 3: Hierarchical question encoding (source: Figure 3 (a) of original paper on page #5)]]<br />
[[File:hierarchy.PNG|thumb|Figure 4: Another figure illustrating hierarchical question encoding in details]]<br />
<br />
=== Word Level ===<br />
1-hot encoding of question's words $Q = \{q_1,..q_T\}$ are transformed into vector space (learned end-to-end) which represents word level embeddings of a visual question i.e. $Q^w = \{q^w_1,...q^w_T\}$. Paper has learned this transformation end-to-end instead of some pretrained models such as word2vec.<br />
<br />
=== Phrase Level ===<br />
Phrase level embedding vectors are calculated by using 1-D convolutions on the word level embedding vectors. <br />
Concretely, at each word location, the inner product of the word vectors with filters of three <br />
window sizes: unigram, bigram and trigram are computed as illustrated by Figure 4. For the ''t-th'' word, <br />
the output from convolution for window size ''s'' is given by<br />
<br />
$$<br />
\hat{q}^p_{s,t} = tanh(W_c^sq^w_{t:t+s-1}), \quad s \in \{1,2,3\}<br />
$$<br />
<br />
Where $W_c^s$ is the weight parameters. The features from three n-grams are combined together using ''maxpool'' operator to obtain the phrase-level embeddings vectors.<br />
<br />
$$<br />
q_t^p = max(\hat{q}^p_{1,t}, \hat{q}^p_{2,t}, \hat{q}^p_{3,t}), \quad t \in \{1,2,...,T\}<br />
$$<br />
<br />
=== Question Level ===<br />
For question level representation, LSTM is used to encode the sequence $q_t^p$ after max-pooling. The corresponding question-level feature at time ''t'' $q_t^s$ is the <br />
LSTM hidden vector at time ''t'' $h_t$.<br />
<br />
$$<br />
\begin{align*}<br />
h_t &= LSTM(q_t^p, h_{t-1})\\<br />
q_t^s &= h_t, \quad t \in \{1,2,...,T\}<br />
\end{align*}<br />
$$<br />
<br />
== Co-Attention Mechanism ==<br />
Paper has proposed two co-attention mechanisms.<br />
{| class="wikitable"<br />
|-<br />
|'''Parallel co-attention'''<br />
|Generates image and question attention simultaneously.<br />
|-<br />
|'''Alternating co-attention'''<br />
|Sequentially alternates between generating image and question attentions.<br />
|}<br />
These co-attention mechanisms are executed at all three levels of the question hierarchy yielding $\hat{v}^r$ and $\hat{q}^r$ <br />
where $r$ is levels in hierarchy i.e. $r \in \{w,p,s\}$ (refer to [[:Notations]] section).<br />
<br />
<br />
=== Parallel Co-Attention ===<br />
[[File:parallewl-coattention.png|thumb|Figure 5: Parallel co-attention mechanism (ref: Figure 2 (a) from original paper)]]<br />
Parallel co-attention attends to the image and question simultaneously as shown in Figure 5. In the paper, "affinity matrix" has been mentioned as the way to calculate the<br />
"attention" or affinity for every pair of image location and question part for each level in the hierarchy (word, phrase and question). Remember, there are $N$ image locations and $T$ <br />
question parts, thus affinity matrix is $\mathbb{R}^{T \times N}$. Specifically, for a given image with<br />
feature map $V \in \mathbb{R}^{d \times N}$, and the question representation $Q \in \mathbb{R}^{d \times T}$, the affinity matrix $C \in \mathbb{R}^{T \times N}$<br />
is calculated by<br />
<br />
$$<br />
C = tanh(Q^TW_bV)<br />
$$<br />
<br />
where,<br />
* $W_b \in \mathbb{R}^{d \times d}$ contains the weights. <br />
<br />
After computing this affinity matrix, one possible way of<br />
computing the image (or question) attention is to simply maximize out the affinity over the locations<br />
of other modality, i.e. $a_v[n] = maxi(C_{i,n})$ and $a_q[t] = maxj(C_{t,j})$. Instead of choosing the max<br />
activation, paper has considered the affinity matrix as a feature and learn to predict image and question attention <br />
maps via the following<br />
<br />
$$<br />
H_v = tanh(W_vV + (W_qQ)C), \quad H_q = tanh(W_qQ + (W_vV )C^T )\\<br />
a_v = softmax(w_{hv}^T Hv), \quad aq = softmax(w_{hq}^T H_q)<br />
$$<br />
<br />
where,<br />
* $W_v, W_q \in \mathbb{R}^{k \times d}$, $w_{hv}, w_{hq} \in \mathbb{R}^k$ are the weight parameters. <br />
* $a_v \in \mathbb{R}^N$ and $a_q \in \mathbb{R}^T$ are the attention probabilities of each image region $v_n$ and word $q_t$ respectively. <br />
<br />
The intuition behind above equation is that, image/question attention maps should be the function of question and image features jointly, therefore, authors have<br />
developed two intermediate parametric functions $H_v$ and $H_q$ that takes affinity matrix $C$, image features $V$ and question features $Q$ as input. The affinity matrix $C$ <br />
transforms question attention space to image attention space (vice versa for $C^T$). Based on the above attention weights, the image and question attention vectors are calculated<br />
as the weighted sum of the image features and question features, i.e.,<br />
<br />
$$\hat{v} = \sum_{n=1}^{N}{a_n^v v_n}, \quad \hat{q}\sum_{t=1}^{T}{a_t^q q_t}$$<br />
<br />
The parallel co-attention is done at each level in the hierarchy, leading to $\hat{v}^r$ and $\hat{q}^r$ where $r \in \{w,p,s\}$. The reason they are using $tanh$ <br />
for $H_q$ and $H_v$is not specified in the paper. But my assumption is that they want to have negative impacts for certain unfavorable pair of image location and question fragment. Unlike $RELU$ or $Sigmoid$, $tanh$ can be between $[-1, 1]$ thus appropriate choice.<br />
<br />
=== Alternating Co-Attention ===<br />
[[File:alternating-coattention.png|thumb|Figure 6: Alternating co-attention mechanism (ref: Figure 2 (b) from original paper)]]<br />
In this attention mechanism, authors sequentially alternate between generating image and question attention as shown in Figure 6. <br />
Briefly, this consists of three steps<br />
<br />
# Summarize the question into a single vector $q$<br />
# Attend to the image based on the question summary $q$<br />
# Attend to the question based on the attended image feature.<br />
<br />
Concretely, paper defines an attention operation $\hat{x} = \mathcal{A}(X, g)$, which takes the image (or question)<br />
features $X$ and attention guidance $g$ derived from question (or image) as inputs, and outputs the<br />
attended image (or question) vector. The operation can be expressed in the following steps<br />
<br />
$$<br />
\begin{align*}<br />
H &= tanh(W_xX + (W_gg)𝟙^T)\\<br />
a_x &= softmax(w_{hx}^T H)\\<br />
\hat{x} &= \sum{a_i^x x_i}<br />
\end{align*}<br />
$$<br />
<br />
where,<br />
* $𝟙$ is a vector with all elements to be 1. <br />
* $W_x, W_g \in \mathbb{R}^{k\times d}$ and $w_{hx} \in \mathbb{R}^k$ are parameters. <br />
* $a_x$ is the attention weight of feature $X$.<br />
<br />
Breifly,<br />
* At the first step of alternating coattention, $X = Q$, and $g$ is $0$. <br />
* At the second step, $X = V$ where $V$ is the image features, and the guidance $g$ is intermediate attended question feature $\hat{s}$ from the first step<br />
* Finally, we use the attended image feature $\hat{v}$ as the guidance to attend the question again, i.e., $X = Q$ and $g = \hat{v}$. <br />
<br />
Similar to the parallel co-attention, the alternating co-attention is also done at each level of the hierarchy, leading to $\hat{v}^r$ <br />
and $\hat{q}^r$ where $r \in \{w,p,s\}$.<br />
<br />
== Encoding for Predicting Answers ==<br />
[[File:answer-encoding-for-prediction.png|thumb|Figure 7: Encoding for predicitng answers (source: Figure 3 (b) of original paper on page #5)]]<br />
Paper treats predicting final answer as a classification task. It was surprising because I always thought answer would a sequence, however, by using MLP it is apparent that answer must be a single word. Co-attended image and question features from all three levels are combined together for the final prediction, see Figure 7. Basically, a multi-layer perceptron (MLP) is deployed to recursively encode the attention features as follows.<br />
$$<br />
\begin{align*}<br />
h_w &= tanh(W_w(\hat{q}^w + \hat{v}^w))\\<br />
h_p &= tanh(W_p[(\hat{q}^p + \hat{v}^p), h_w])\\<br />
h_s &= tanh(W_s[(\hat{q}^s + \hat{v}^s), h_p])\\<br />
p &= softmax(W_hh^s)<br />
\end{align*}<br />
$$<br />
<br />
where <br />
* $W_w, W_p, W_s$ and $W_h$ are the weight parameters. <br />
* $[·]$ is the concatenation operation on two vectors. <br />
* $p$ is the probability of the final answer.<br />
<br />
= Experiments =<br />
Evaluation for the proposed model is performed using two datasets, the VQA dataset [1] and the COCO-QA dataset [2].<br />
<br />
* '''VQA dataset''' is the largest dataset for this problem, containing human annotated questions and answers on Microsoft COCO dataset.<br />
* '''COCO-QA dataset''' is automatically generated from captions in the Microsoft COCO dataset.<br />
<br />
The proposed approach seems to outperform most of the state-of-art techniques as shown in Table 1 and 2.<br />
<br />
[[File:result-vqa.png|thumb|700px|center|Table 1: Results on the VQA dataset. “-” indicates the results is not available. (ref: Table 1 of original paper page #6)]]<br />
<br />
[[File:result-coco-qa.png|thumb|700px|center|Table 2: Results on the COCO-QA dataset. “-” indicates the results is not available (ref: Table 2 of original paper page #7)]]<br />
<br />
= Qualitative Results =<br />
We now visualize some co-attention maps generated by their method in Figure 8. <br />
<br />
{|class="wikitable"<br />
|'''Word level'''<br />
|<br />
* Model attends mostly to the object regions in an image, and objects at questions as well e.g., heads, bird. <br />
|-<br />
|'''Phrase level'''<br />
|<br />
*Image attention has different patterns across images. <br />
** For the first two images, the attention transfers from objects to background regions. <br />
** For the third image, the attention becomes more focused on the objects. <br />
** Reason for different attention could be perhaphs caused by the different question types. <br />
* On the question side, their model is capable of localizing the key phrases in the question, thus essentially discovering the question types in the dataset. <br />
* For example, their model pays attention to the phrases “what color” and “how many snowboarders”. <br />
|-<br />
|'''Question level'''<br />
|<br />
* Image attention concentrates mostly on objects. <br />
* Their model successfully attends to the regions in images and phrases in the questions appropriate for answering the question, e.g., “color of the bird” and bird region.<br />
|}<br />
<br />
Because their model performs co-attention at three levels, it often captures complementary information from<br />
each level, and then combines them to predict the answer. However, it some what un-intuitive to visualize the <br />
phrase and question level attention mapping applied diretly to the words of the question, since phrase <br />
and question level features are compound features from multiple words, thus their attention contribution on the <br />
actual words from the question cannot be clearly understood. <br />
<br />
[[File:visualization-co-attention.png|thumb|800px|center|Figure 8: Visualization of image and question co-attention maps on the COCO-QA dataset. From left to right:<br />
original image and question pairs, word level co-attention maps, phrase level co-attention maps and question<br />
level co-attention maps. For visualization, both image and question attentions are scaled (from red:high to<br />
blue:low). (ref: Figure 4 of original paper page #8)]]<br />
<br />
= Conclusion =<br />
* A hierarchical co-attention model for visual question answering is proposed. <br />
* Coattention allows model to attend to different regions of the image as well as different fragments of the question. <br />
* Question is hierarchically represented at three levels to capture information from different granularities. <br />
* Visualization shows model co-attends to interpretable regions of images and questions for predicting the answer. <br />
* Though their model was evaluated on visual question answering, it can be potentially applied to other tasks involving vision and language.<br />
<br />
= Reference =<br />
# K. Kafle and C. Kanan, “Visual Question Answering: Datasets, Algorithms, and Future Challenges,” Computer Vision and Image Understanding, Jun. 2017.<br />
# Mengye Ren, Ryan Kiros, and Richard Zemel. Exploring models and data for image question answering. NIPS, 2015.</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Hierarchical_Question-Image_Co-Attention_for_Visual_Question_Answering&diff=31292Hierarchical Question-Image Co-Attention for Visual Question Answering2017-11-23T18:08:33Z<p>Nkadiyar: </p>
<hr />
<div>__TOC__<br />
== Paper Summary ==<br />
{| class="wikitable"<br />
|-<br />
|'''Conference'''<br />
| <br />
* NIPS 2016<br />
* Presented as spotlight oral: [https://www.youtube.com/watch?v=m6t9IFdk0ms Youtube link]<br />
* 85 citations so far<br />
|-<br />
| '''Authors'''<br />
|Jiasen Lu, Jianwei Yang, Dhruv Batra, '''Devi Parikh'''<br />
|-<br />
|'''Abstract'''<br />
|''A number of recent works have proposed attention models for Visual Question Answering (VQA) that generate spatial maps highlighting image regions relevant to answering the question. In this paper, we argue that in addition to modeling "where to look" or visual attention, it is equally important to model "what words to listen to" or question attention. We present a novel co-attention model for VQA that jointly reasons about image and question attention. In addition, our model reasons about the question (and consequently the image via the co-attention mechanism) in a hierarchical fashion via a novel 1-dimensional convolution neural networks (CNN). Our model improves the state-of-the-art on the VQA dataset from 60.3% to 60.5%, and from 61.6% to 63.3% on the COCO-QA dataset. By using ResNet, the performance is further improved to 62.1% for VQA and 65.4% for COCO-QA.''<br />
|}<br />
= Introduction =<br />
'''Visual Question Answering (VQA)''' is a recent problem in computer vision and<br />
natural language processing that has garnered a large amount of interest from<br />
the deep learning, computer vision, and natural language processing communities.<br />
In VQA, an algorithm needs to answer text-based questions about images in<br />
natural language as illustrated in Figure 1.<br />
<br />
[[File:vqa-overview.png|thumb|600px|center|Figure 1: Illustration of VQA system whereby machine learning algorithm answers a visual question asked by an user for a given image (ref: http://www.visualqa.org/static/img/challenge.png)]]<br />
<br />
Recently, ''visual-attention'' based models have gained traction for VQA tasks, where the<br />
attention mechanism typically produces a spatial map highlighting image regions<br />
relevant for answering the visual question about the image. However, to correctly answer the <br />
question, machine not only needs to understand or "attend"<br />
regions in the image but also the parts of question as well. In this paper, authors have proposed a novel ''co-attention''<br />
technique to combine "where to look" or visual-attention along with "what words<br />
to listen to" or question-attention VQA allowing their model to jointly reasons about image and question thus improving <br />
upon existing state of the art results.<br />
<br />
== "Attention" Models ==<br />
You may skip this section if you already know about "attention" in<br />
context of deep learning. Since this paper talks about "attention" almost<br />
everywhere, I decided to put this section to give very informal and brief<br />
introduction to the concept of the "attention" mechanism specially visual "attention", <br />
however, it can be expanded to any other type of "attention".<br />
<br />
Visual attention in CNN is inspired by the biological visual system. As humans,<br />
we have ability to focus our cognitive processing onto a subset of the<br />
environment that is more relevant for the given situation. Imagine, you witness<br />
a bank robbery where robbers are trying to escape on a car, as a good citizen,<br />
you will immediately focus your attention on number plate and other physical<br />
features of the car and robbers in order to give your testimony later, however, you may not remember things which otherwise interests you more. <br />
Such selective visual attention for a given context (robbery in above example) can also be implemented in<br />
traditional CNNs as well. This allows CNNs to be more robust and superior for certain tasks and it even helps <br />
algorithm designer to visualize what spacial features (regions within image) were more important than others. Attention guided<br />
deep learning is particularly very helpful for image caption and VQA tasks.<br />
<br />
== Role of Visual Attention in VQA ==<br />
This section is not a part of the actual paper that is been summarized, however, it gives an overview<br />
of how visual attention can be incorporated in training of a network for VQA tasks, eventually, helping <br />
readers to absorb and understand actual proposed ideas from the paper more effortlessly.<br />
<br />
Generally for implementing attention, network tries to learn the conditional <br />
distribution $P_{i \in [1,n]}(Li|c)$ representing individual importance for all the features <br />
extracted from each of the dsicrete $n$ locations within the image <br />
conditioned on some context vector $c$. In order words, given $n$ features <br />
$L_i = [L_1, ..., L_n]$ from $n$ different spacial regions within the image (top-left, top-middle, top-right, and so on), <br />
then "attention" module learns a parameteric function $F(c;\theta)$ that outputs an importance mapping <br />
of each of these individual feature for a given context vector $c$ or a discrete probability distribution <br />
of size $n$, can be achived by $softmax(n)$. <br />
<br />
In order to incorporate the visual attention in VQA task, one can define context vector $c$ <br />
as a representation of the visual question asked by an user (using RNN perhaps LSTM). The context $c$ can then be used to generate an <br />
attention map for corresponding image locations (as shown in Figure 2) further improving the accuracy on final end-to-end training. <br />
Most work that exists in literature regarding use of visual-attention in VQA tasks are generally further <br />
specialization of the similar ideas.<br />
<br />
[[File:attention-vqa-general.png|thumb|700px|center|Figure 2: Different attention maps generated based on the given visual question. Regions with most "attention" or importance is whitened, machine learning model has learned to steer its attention based on the given question.]]<br />
<br />
== Motivation and Main Contributions ==<br />
So far, all attention models for VQA in literature have focused on the problem of identifying "where<br />
to look" or visual attention. In this paper, authors argue that the problem of identifying "which words to<br />
listen to" or '''question attention''' is equally important. Consider the questions "how many horses are<br />
in this image?" and "how many horses can you see in this image?". They have the same meaning,<br />
essentially captured by the first three words. A machine that attends to the first three words would<br />
arguably be more robust to linguistic variations irrelevant to the meaning and answer of the question.<br />
Motivated by this observation, in addition to reasoning about visual attention, paper has addressed the<br />
problem of question attention. Basically, main contributions of the paper are as follows.<br />
<br />
* A novel co-attention mechanism for VQA that jointly performs question-guided visual attention and image-guided question attention.<br />
* A hierarchical architecture to represent the question, and consequently construct image-question co-attention maps at 3 different levels: word level, phrase level and question level. <br />
* A novel convolution-pooling strategy at phase-level to adaptively select the phrase sizes whose representations are passed to the question level representation.<br />
* Results on VQA and COCO-QA and ablation studies to quantify the roles of different components in the model<br />
<br />
= Method =<br />
This section is broken down into four parts: '''(i)''' notations used within the paper and also throughout this summary, '''(ii)''' hierarchical representation for a visual question, '''(iii)''' the proposed co-attention mechanism and<br />
'''(iv)''' predicting answers.<br />
<br />
== Notations ==<br />
{| class="wikitable"<br />
|-<br />
|'''Notation'''<br />
|'''Explaination'''<br />
|-<br />
|$Q = \{q_1,...q_T\}$<br />
|One-hot encoding of a visual question with $T$ words. Paper uses three different representation of visual question, one for each level of hierarchy, they are as follows: <br />
# $Q^w = \{q^w_1,...q^w_T\}$: Word level representation of visual question<br />
# $Q^p = \{q^p_1,...q^p_T\}$: Phrase level representation of visual question<br />
# $Q^s = \{q^s_1,...q^s_T\}$: Question level representation of visual question<br />
$Q^{w,p,s}$ has exactly $T$ number of embeddings in it (sequential data with temporal dimension), regardless of its position in the hierarchy i.e. word, phrase or question. <br />
|-<br />
|$V = {v_1,..,v_N}$<br />
|$V$ represented various vectors from $N$ different locations within the given image. Therefore, $v_n$ is feature vector from the image at location $n$. $V$ collectively covers the entire spatial reachings of the image. One can extract these location sensitive features from convolution layer of CNN.<br />
|-<br />
|$\hat{v}^r$ and $\hat{q}^r$<br />
|The co-attention features of image and question at each level in the hierarchy where $r \in \{w,p,s\}$. Basically, its a sum of $Q$ or $V$ after the dot product with attention $a^q$ or $a^v$ at each level of hierarchy. <br />
For example, at "word" level, $a^q_w$ and $a^v_w$ is a probablity distribution representing importance of each words in visual question and each location within image respectively, whereas $\hat{q}^w$ and $\hat{v}^w$ are final features vectors for the given question and image with attention maps ($a^q_w$ and $a^v_w$ applied) at the "word" level, and similary for "phrase" and "question" level as well.<br />
|}<br />
'''Note:''' Throughout the paper, $W$ represents the learnable weights and biases are not used within the equations for simplicity (reader must assume it to exist).<br />
<br />
== Question Hierarchy ==<br />
There are three levels of granularities for their hierarchical representation of a visual question: '''(i)''' word, '''(ii)''' phrase and '''(iii)''' question level. It is important to note, each level depends on the previous one, so, phrase level representations are extracted from word level and question level representations come from phrase level as depicted in Figure 4.<br />
<br />
[[File:hierarchy2.png|thumb|Figure 3: Hierarchical question encoding (source: Figure 3 (a) of original paper on page #5)]]<br />
[[File:hierarchy.PNG|thumb|Figure 4: Another figure illustrating hierarchical question encoding in details]]<br />
<br />
=== Word Level ===<br />
1-hot encoding of question's words $Q = \{q_1,..q_T\}$ are transformed into vector space (learned end-to-end) which represents word level embeddings of a visual question i.e. $Q^w = \{q^w_1,...q^w_T\}$. Paper has learned this transformation end-to-end instead of some pretrained models such as word2vec.<br />
<br />
=== Phrase Level ===<br />
Phrase level embedding vectors are calculated by using 1-D convolutions on the word level embedding vectors. <br />
Concretely, at each word location, the inner product of the word vectors with filters of three <br />
window sizes: unigram, bigram and trigram are computed as illustrated by Figure 4. For the ''t-th'' word, <br />
the output from convolution for window size ''s'' is given by<br />
<br />
$$<br />
\hat{q}^p_{s,t} = tanh(W_c^sq^w_{t:t+s-1}), \quad s \in \{1,2,3\}<br />
$$<br />
<br />
Where $W_c^s$ is the weight parameters. The features from three n-grams are combined together using ''maxpool'' operator to obtain the phrase-level embeddings vectors.<br />
<br />
$$<br />
q_t^p = max(\hat{q}^p_{1,t}, \hat{q}^p_{2,t}, \hat{q}^p_{3,t}), \quad t \in \{1,2,...,T\}<br />
$$<br />
<br />
=== Question Level ===<br />
For question level representation, LSTM is used to encode the sequence $q_t^p$ after max-pooling. The corresponding question-level feature at time ''t'' $q_t^s$ is the <br />
LSTM hidden vector at time ''t'' $h_t$.<br />
<br />
$$<br />
\begin{align*}<br />
h_t &= LSTM(q_t^p, h_{t-1})\\<br />
q_t^s &= h_t, \quad t \in \{1,2,...,T\}<br />
\end{align*}<br />
$$<br />
<br />
== Co-Attention Mechanism ==<br />
Paper has proposed two co-attention mechanisms.<br />
{| class="wikitable"<br />
|-<br />
|'''Parallel co-attention'''<br />
|Generates image and question attention simultaneously.<br />
|-<br />
|'''Alternating co-attention'''<br />
|Sequentially alternates between generating image and question attentions.<br />
|}<br />
These co-attention mechanisms are executed at all three levels of the question hierarchy yielding $\hat{v}^r$ and $\hat{q}^r$ <br />
where $r$ is levels in hierarchy i.e. $r \in \{w,p,s\}$ (refer to [[:Notations]] section).<br />
<br />
<br />
=== Parallel Co-Attention ===<br />
[[File:parallewl-coattention.png|thumb|Figure 5: Parallel co-attention mechanism (ref: Figure 2 (a) from original paper)]]<br />
Parallel co-attention attends to the image and question simultaneously as shown in Figure 5. In the paper, "affinity matrix" has been mentioned as the way to calculate the<br />
"attention" or affinity for every pair of image location and question part for each level in the hierarchy (word, phrase and question). Remember, there are $N$ image locations and $T$ <br />
question parts, thus affinity matrix is $\mathbb{R}^{T \times N}$. Specifically, for a given image with<br />
feature map $V \in \mathbb{R}^{d \times N}$, and the question representation $Q \in \mathbb{R}^{d \times T}$, the affinity matrix $C \in \mathbb{R}^{T \times N}$<br />
is calculated by<br />
<br />
$$<br />
C = tanh(Q^TW_bV)<br />
$$<br />
<br />
where,<br />
* $W_b \in \mathbb{R}^{d \times d}$ contains the weights. <br />
<br />
After computing this affinity matrix, one possible way of<br />
computing the image (or question) attention is to simply maximize out the affinity over the locations<br />
of other modality, i.e. $a_v[n] = maxi(C_{i,n})$ and $a_q[t] = maxj(C_{t,j})$. Instead of choosing the max<br />
activation, paper has considered the affinity matrix as a feature and learn to predict image and question attention <br />
maps via the following<br />
<br />
$$<br />
H_v = tanh(W_vV + (W_qQ)C), \quad H_q = tanh(W_qQ + (W_vV )C^T )\\<br />
a_v = softmax(w_{hv}^T Hv), \quad aq = softmax(w_{hq}^T H_q)<br />
$$<br />
<br />
where,<br />
* $W_v, W_q \in \mathbb{R}^{k \times d}$, $w_{hv}, w_{hq} \in \mathbb{R}^k$ are the weight parameters. <br />
* $a_v \in \mathbb{R}^N$ and $a_q \in \mathbb{R}^T$ are the attention probabilities of each image region $v_n$ and word $q_t$ respectively. <br />
<br />
The intuition behind above equation is that, image/question attention maps should be the function of question and image features jointly, therefore, authors have<br />
developed two intermediate parameteric functions $H_v$ and $H_q$ that takes affinity matrix $C$, image features $V$ and question features $Q$ as input. The affinity matrix $C$ <br />
transforms question attention space to image attention space (vice versa for $C^T$). Based on the above attention weights, the image and question attention vectors are calculated<br />
as the weighted sum of the image features and question features, i.e.,<br />
<br />
$$\hat{v} = \sum_{n=1}^{N}{a_n^v v_n}, \quad \hat{q}\sum_{t=1}^{T}{a_t^q q_t}$$<br />
<br />
The parallel co-attention is done at each level in the hierarchy, leading to $\hat{v}^r$ and $\hat{q}^r$ where $r \in \{w,p,s\}$. The reason they are using $tanh$ <br />
for $H_q$ and $H_v$is not specified in the paper. But my assumption is that they want to have negative imapcts for certain unfavorable pair of image location and question fragment. Unlike $RELU$ or $Sigmoid$, $tanh$ can be betwen $[-1, 1]$ thus appropriate choice.<br />
<br />
=== Alternating Co-Attention ===<br />
[[File:alternating-coattention.png|thumb|Figure 6: Alternating co-attention mechanism (ref: Figure 2 (b) from original paper)]]<br />
In this attention mechanism, authors sequentially alternate between generating image and question attention as shown in Figure 6. <br />
Briefly, this consists of three steps<br />
<br />
# Summarize the question into a single vector $q$<br />
# Attend to the image based on the question summary $q$<br />
# Attend to the question based on the attended image feature.<br />
<br />
Concretely, paper defines an attention operation $\hat{x} = \mathcal{A}(X, g)$, which takes the image (or question)<br />
features $X$ and attention guidance $g$ derived from question (or image) as inputs, and outputs the<br />
attended image (or question) vector. The operation can be expressed in the following steps<br />
<br />
$$<br />
\begin{align*}<br />
H &= tanh(W_xX + (W_gg)𝟙^T)\\<br />
a_x &= softmax(w_{hx}^T H)\\<br />
\hat{x} &= \sum{a_i^x x_i}<br />
\end{align*}<br />
$$<br />
<br />
where,<br />
* $𝟙$ is a vector with all elements to be 1. <br />
* $W_x, W_g \in \mathbb{R}^{k\times d}$ and $w_{hx} \in \mathbb{R}^k$ are parameters. <br />
* $a_x$ is the attention weight of feature $X$.<br />
<br />
Breifly,<br />
* At the first step of alternating coattention, $X = Q$, and $g$ is $0$. <br />
* At the second step, $X = V$ where $V$ is the image features, and the guidance $g$ is intermediate attended question feature $\hat{s}$ from the first step<br />
* Finally, we use the attended image feature $\hat{v}$ as the guidance to attend the question again, i.e., $X = Q$ and $g = \hat{v}$. <br />
<br />
Similar to the parallel co-attention, the alternating co-attention is also done at each level of the hierarchy, leading to $\hat{v}^r$ <br />
and $\hat{q}^r$ where $r \in \{w,p,s\}$.<br />
<br />
== Encoding for Predicting Answers ==<br />
[[File:answer-encoding-for-prediction.png|thumb|Figure 7: Encoding for predicitng answers (source: Figure 3 (b) of original paper on page #5)]]<br />
Paper treats predicitng final answer as a classification task. It was surprising because I always thought answer would a sequence, however, by using MLP it is apparent that answer must be a single word. Co-attended image and question features from all three levels are combined together for the final prediction, see Figure 7. Basically, a multi-layer perceptron (MLP) is deployed to recursively encode the attention features as follows.<br />
$$<br />
\begin{align*}<br />
h_w &= tanh(W_w(\hat{q}^w + \hat{v}^w))\\<br />
h_p &= tanh(W_p[(\hat{q}^p + \hat{v}^p), h_w])\\<br />
h_s &= tanh(W_s[(\hat{q}^s + \hat{v}^s), h_p])\\<br />
p &= softmax(W_hh^s)<br />
\end{align*}<br />
$$<br />
<br />
where <br />
* $W_w, W_p, W_s$ and $W_h$ are the weight parameters. <br />
* $[·]$ is the concatenation operation on two vectors. <br />
* $p$ is the probability of the final answer.<br />
<br />
= Experiments =<br />
Evaluation for the proposed model is performed using two datasets, the VQA dataset [1] and the COCO-QA dataset [2].<br />
<br />
* '''VQA dataset''' is the largest dataset for this problem, containing human annotated questions and answers on Microsoft COCO dataset.<br />
* '''COCO-QA dataset''' is automatically generated from captions in the Microsoft COCO dataset.<br />
<br />
The proposed approach seems to outperform most of the state-of-art techniques as shown in Table 1 and 2.<br />
<br />
[[File:result-vqa.png|thumb|700px|center|Table 1: Results on the VQA dataset. “-” indicates the results is not available. (ref: Table 1 of original paper page #6)]]<br />
<br />
[[File:result-coco-qa.png|thumb|700px|center|Table 2: Results on the COCO-QA dataset. “-” indicates the results is not available (ref: Table 2 of original paper page #7)]]<br />
<br />
= Qualitative Results =<br />
We now visualize some co-attention maps generated by their method in Figure 8. <br />
<br />
{|class="wikitable"<br />
|'''Word level'''<br />
|<br />
* Model attends mostly to the object regions in an image, and objects at questions as well e.g., heads, bird. <br />
|-<br />
|'''Phrase level'''<br />
|<br />
*Image attention has different patterns across images. <br />
** For the first two images, the attention transfers from objects to background regions. <br />
** For the third image, the attention becomes more focused on the objects. <br />
** Reason for different attention could be perhaphs caused by the different question types. <br />
* On the question side, their model is capable of localizing the key phrases in the question, thus essentially discovering the question types in the dataset. <br />
* For example, their model pays attention to the phrases “what color” and “how many snowboarders”. <br />
|-<br />
|'''Question level'''<br />
|<br />
* Image attention concentrates mostly on objects. <br />
* Their model successfully attends to the regions in images and phrases in the questions appropriate for answering the question, e.g., “color of the bird” and bird region.<br />
|}<br />
<br />
Because their model performs co-attention at three levels, it often captures complementary information from<br />
each level, and then combines them to predict the answer. However, it some what un-intuitive to visualize the <br />
phrase and question level attention mapping applied diretly to the words of the question, since phrase <br />
and question level features are compund features from multiple words, thus their attention contribution on the <br />
actual words from the question cannot be clearly understood. <br />
<br />
[[File:visualization-co-attention.png|thumb|800px|center|Figure 8: Visualization of image and question co-attention maps on the COCO-QA dataset. From left to right:<br />
original image and question pairs, word level co-attention maps, phrase level co-attention maps and question<br />
level co-attention maps. For visualization, both image and question attentions are scaled (from red:high to<br />
blue:low). (ref: Figure 4 of original paper page #8)]]<br />
<br />
= Conclusion =<br />
* A hierarchical co-attention model for visual question answering is proposed. <br />
* Coattention allows model to attend to different regions of the image as well as different fragments of the question. <br />
* Question is hierarchically represented at three levels to capture information from different granularities. <br />
* Visualization shows model co-attends to interpretable regions of images and questions for predicting the answer. <br />
* Though their model was evaluated on visual question answering, it can be potentially applied to other tasks involving vision and language.<br />
<br />
= Reference =<br />
# K. Kafle and C. Kanan, “Visual Question Answering: Datasets, Algorithms, and Future Challenges,” Computer Vision and Image Understanding, Jun. 2017.<br />
# Mengye Ren, Ryan Kiros, and Richard Zemel. Exploring models and data for image question answering. NIPS, 2015.</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Conditional_Image_Synthesis_with_Auxiliary_Classifier_GANs&diff=31291Conditional Image Synthesis with Auxiliary Classifier GANs2017-11-23T18:07:08Z<p>Nkadiyar: </p>
<hr />
<div>'''Abstract:''' "In this paper we introduce new methods for the improved training of generative adversarial networks (GANs) for image synthesis. We construct a variant of GANs employing label conditioning that results in 128×128 resolution image samples exhibiting global coherence. We expand on previous work for image quality assessment to provide two new analyses for assessing the discriminability and diversity of samples from class-conditional image synthesis models. These analyses demonstrate that high resolution samples provide class information not present in low resolution samples. Across 1000 ImageNet classes, 128×128 samples are more than twice as discriminable as artificially resized 32×32 samples. In addition, 84.7% of the classes have samples exhibiting diversity comparable to real ImageNet data." [[#References | (Odena et al., 2016)]]<br />
<br />
= Introduction =<br />
<br />
=== Motivation ===<br />
The authors introduce a GAN architecture for generating high resolution images from the ImageNet dataset. They show that this architecture makes it possible to split the generation process into many sub-models. They further suggest that GANs have trouble generating globally coherent images, and that this architecture is responsible for the coherence of their samples. They experimentally demonstrate that generating higher resolution images allow the model to encode more class-specific information, making them more visually discriminable than lower resolution images even after they have been resized to the same resolution.<br />
<br />
The second half of the paper introduces metrics for assessing visual discriminability and diversity of synthesized images. The discussion of image diversity in particular is important due to the tendency for GANs to 'collapse' to only produce one image that best fools the discriminator [[#References|(Goodfellow et al., 2014)]].<br />
<br />
=== Previous Work ===<br />
<br />
Of all image synthesis methods (e.g. variational autoencoders, autoregressive models, invertible density estimators), GANs have become one of the most popular and successful due to their flexibility and the ease with which they can be sampled from. A standard GAN framework pits a generative model $G$ against a discriminative adversary $D$. The goal of $G$ is to learn a mapping from a latent space $Z$ to a real space $X$ to produce examples (generally images) indistinguishable from training data. The goal of the $D$ is to iteratively learn to predict when a given input image is from the training set or a synthesized image from $G$. Jointly the models are trained to solve the game-theoretical minimax problem, as defined by [[#References|Goodfellow et al. (2014)]]: $$\underset{G}{\text{min }}\underset{D}{\text{max }}V(G,D)=\mathbb{E}_{X\sim p_{data}(x)}[log(D(X))]+\mathbb{E}_{Z\sim p_{Z}(z)}[log(1-D(G(Z)))]$$<br />
<br />
While this initial framework has clearly demonstrated great potential, other authors have proposed changes to the method to improve it. Many such papers propose changes to the training process [[#References|(Salimans et al., 2016)]][[#References|(Karras et al., 2017)]], which is notoriously difficult for some problems. Others propose changes to the model itself. [[#References|Mirza & Osindero (2014)]] augment the model by supplying the class of observations to both the generator and discriminator to produce class-conditional samples. According to [[STAT946F17/Conditional Image Generation with PixelCNN Decoders|van den Oord et al. (2016)]], conditioning image generation on classes can greatly improve their quality. Other authors have explored using even richer side information in the generation process with good results [[Learning What and Where to Draw|(Reed et al., 2016)]].<br />
<br />
Another model modification relevant to this paper is to force the discriminator network to reconstruct side information by adding an auxiliary network to classify generated (and real) images. The authors make the claim that forcing a model to perform additional tasks is known to improve performance on the original task [[#References|(Szegedy et al., 2014)]][[#References|(Sutskever et al., 2014)]][[#References|(Ramsundar et al., 2016)]]. They further suggest that using pre-trained image classifiers (rather than classifiers trained on both real and generated images) could improve results over and above what is shown in this paper.<br />
<br />
= Contributions =<br />
<br />
The contributions of this paper are in three main areas. First, the authors propose slight changes to previously existing GAN architectures, resulting in a model capable of generating samples of impressive quality. Second, the authors propose two metrics to assess the quality of samples generated from a GAN. Lastly, they present empirical results on GANs which are of some interest.<br />
<br />
== Model ==<br />
<br />
The authors propose an auxiliary classifier GAN (AC-GAN) which is a slight variation on previous architectures. Like [[#References|Mirza & Osindero (2014)]], the generator takes the image class to be generated as input in addition to the latent encoding $Z$. Like [[#References|Odena (2016)]] and [[#References|Salimans et al. (2016)]], the discriminator is trained to predict not only whether an observation is real or fake, but to classify each observation as well. The marginal contribution of this paper is to combine these in one model.<br />
<br />
Formally, let $C\sim p_c$ represent the target class label of each generated observation and $Z$ represent the usual noise vector from the latent space. Then the generator function takes both as argument to produce image samples: $X_{fake}=G(c,z)$.The discriminator gives a probability distribution over the source $S$ (real or fake) of the image as well as the class label $C$ being generated. $$D(X)=<P(S|X),P(C|X)>$$<br />
<br />
The objective function for the model thus has two parts, one corresponding to the source $L_S$ and the other to the class $L_C$. $D$ is trained to maximize $L_S + L_C$, while $G$ is trained to maximize $L_C-L_S$. Using the notation of Goodfellow et al. (2014), the loss terms are defined as:<br />
$$L_S=\mathbb{E}_{X\sim p_{data}(x)}[log(P(S=real|X))]+\mathbb{E}_{C,Z\sim p_{C,Z}(c,z)}[log(P(S=fake|G(C,Z)))]$$<br />
$$L_C=\mathbb{E}_{X\sim p_{data}(x)}[log(P(C|X))]+\mathbb{E}_{C,Z\sim p_{C,Z}(c,z)}[log(P(C|G(C,Z)))]$$<br />
<br />
Because G accepts both $C$ and $Z$ as arguments, it is able to learn a mapping $Z\rightarrow X$ that is independent of $C$. The authors argue that all class-specific information should be represented by $C$, allowing $Z$ to represent other factors such as pose, background, etc.<br />
<br />
Lastly the authors split the generation process into many class-specific submodels. They point out that the structure of their model permits this split, though it should technically be possible for even the standard GAN framework by dividing the training data into groups according to their known class labels.<br />
<br />
The changes above result in a model capable of generating (some) image samples with both high resolution and global coherence.<br />
<br />
== GAN Quality Metrics ==<br />
<br />
A much larger part of the paper is spent on measuring the quality of a GAN's output. As the authors say, evaluating a generative model's quality is difficult due to a large number of probabilistic measures (such as average log-likelihood, Parzen window estimates, and visual fidelity [[#References| (Theis et al., 2015)]]) and "a lack of a perceptually meaningful image similarity metric".<br />
<br />
=== Image Discriminability Metric ===<br />
The authors develop two metrics in this paper to address these shortcomings. The first of these is a discriminability metric, the goal of which is to assess the degree to which generated images are identifiable as the class they are meant to represent. Ideally a team of non-expert humans could handle this, but the difficulty of such an approach makes the need for an automated metric apparent. The metric proposed by the authors is to measure the accuracy of a pre-trained image classifier trained on the pristine training data. For this they select a [https://github.com/openai/improved-gan/ modified version of Inception-v3][[#References|(Szegedy et al., 2015)]].<br />
<br />
Other metrics already exist for assessing image quality, the most popular of which is probably the Inception Score [[#References| (Salimans et al., 2016)]]. The authors list two main advantages of their approach over the Inception Score. The first is that accuracy figures are easier to interpret than the Inception Score, which is fairly self-evident. The second advantage of using Inception accuracy instead of Inception Score is that Inception accuracy may be calculated for individual classes, giving a better picture of where the model is strong and where it is weak.<br />
<br />
=== Image Diversity Metric ===<br />
<br />
The second metric proposed by the authors measures the diversity of the generated images. As mentioned above, image diversity is an important quality in a GAN, a common failure they suffer from is 'collapsing', where the generator learns to only output one image that is good at fooling the discriminator [[#References|(Goodfellow et al., 2014)]][[#References|(Salimans et al., 2016)]]. The metric proposed in this section is intended to be complementary to the Inception accuracy, as Inception accuracy would not detect generator collapse.<br />
<br />
For their diversity metric, the authors co-opt an existing metric used to measure the similarity between two images: multi-scale structural similarity (MS-SSIM)[[#References|(Wang et al., 2003)]]. The authors do not go into detail about how the MS-SSIM measure is calculated, except to say that it is one of the more successful ways to predict human perceptual similarity judgements. It can take values on the interval $[0,1]$, and higher values indicate the two images being compared a perceptually more similar. For images generated from a GAN then, the metric should ideally be low, as diversity is desired.<br />
<br />
The authors contribution is to use this metric to assess diversity of a GAN's output. It is a pairwise comparison between two images, so their solution is to compare 100 images (that is $100\cdot 99$ paired comparisons) from each class and take the mean MS-SSIM score.<br />
<br />
The authors make two points about their use of this metric. First, the way they apply the metric is different from how it was originally intended to be used. It is possible that it will not behave as desired because of this. As evidence to the contrary, they state that:<br />
# Visually the metric seems to work. Pairs with high MS-SSIM seem more similar.<br />
# Comparisons are only made between images in the same class, keeping their application of the metric closer to its original use case of measuring the quality of compression algorithms.<br />
# The metric is not saturated. Scores on their generated data vary across the unit interval. If scores were all very close to zero the metric would not be much use.<br />
<br />
The second point they raise is that the mean MS-SSIM metric is not intended as a proxy for entropy of the generator distribution in pixel space. That measure is hard to compute, and in any case is sensitive to trivial changes in the pixels, whereas the true intention of this metric is to measure perceptual similarity.<br />
<br />
== Experimental Results on GAN Properties ==<br />
<br />
[[File:Figure_2_(Bottom).JPG|thumb|500px|right|alt=(Odena et al., 2016) Figure 2: (Left) Inception accuracy (y-axis) of two generators with resolution 128 x 128 (red) and 64 x 64 (blue). Images are resized to the same spatial resolution (x-axis). (Right) Class accuracies from the 128 x 128 AC-GAN at full resolution (x-axis) and downsized to 32 x 32 (y-axis).|(Odena et al., 2016) Figure 2: (Left) Inception accuracy (y-axis) of two generators with resolution 128 x 128 (red) and 64 x 64 (blue). Images are resized to the same spatial resolution (x-axis). (Right) Class accuracies from the 128 x 128 AC-GAN at full resolution (x-axis) and downsized to 32 x 32 (y-axis).]]<br />
<br />
The authors conduct several experiments on their model and proposed metrics. These are summarized in this section.<br />
<br />
=== Higher Resolution Images are more Discriminable ===<br />
<br />
As it is one of the main attractions of this paper, the authors investigate how generating samples at different resolutions affects their discriminability. To achieve this two models are trained, one that generates 64 x 64 resolution images and one that generates 128 x 128 resolution images. These images can be rescaled using bilinear interpolation to make them directly comparable. The authors find that the 128 x 128 AC-GAN (on average) achieves higher discriminability (per the Inception accuracy metric [[#Image Discriminability Metric|introduced above]]) at all resized resolutions. The authors claim that theirs is the first work to investigate how much an image generator is 'making use of its given output resolution'.<br />
<br />
=== Generated Images are both Diverse and Discriminable ===<br />
<br />
A second experiment conducted in the paper aims to investigate how the two metrics they propose interact with each other. They simply calculate the Inception accuracy and mean MS-SSIM score for a batch of generated images from every class in their data and report the correlation between the scores. They find the scores are anti-correlated with $\rho=-0.16$. Because the mean MS-SSIM metric is low for diverse samples, they conclude that accuracy and diversity are actually positively correlated, which is contradictory to the hypothesis that GANs that collapse achieve better sample quality.<br />
<br />
[[File:Figure_10.JPG|thumb|300px|right|alt=(Odena et al., 2016) Figure 10: Mean MS-SSIM scores for 10 ImageNet classes (y-axis) plotted against the number of classes handled by each sub-model (x-axis).|(Odena et al., 2016) Figure 10: Mean MS-SSIM scores for 10 ImageNet classes (y-axis) plotted against the number of classes handled by each sub-model (x-axis).]]<br />
<br />
=== Effect of Class Splits on Image Sample Quality ===<br />
<br />
The final experiment conducted in the paper investigates how the number of classes a fixed GAN model has to generate impacts the diversity of the generated images. In their main model the authors split the data such that each AC-GAN only had to generate images from 10 classes. Here they experimented with changing the number of classes each sub-model has to generate while holding the architecture of the sub-models fixed. They report the mean MS-SSIM score of the generated images from the original classes (when the model had only 10 to generate) at each split level. Perhaps unsurprisingly, the diversity of the outputs drops as the number of classes the model has to handle increases. Giving the model more parameters to handle the larger (and more diverse) set of classes might possibly eliminate this effect.<br />
<br />
Finally, they state that they were unable to get a regular GAN to converge on the task of generating images from one class per model. This could be due to the difficulty of training GANs and the limited amount of training data available for each class rather than any theoretical property of class splits.<br />
<br />
= Results =<br />
<br />
=== Model ===<br />
The authors apply their AC-GAN model to the ImageNet [[#References|(Russakovsky et al., 2015)]] dataset. The data have 1000 classes which the authors split into 100 groups of 10. An AC-GAN model is trained on each group of 10 to give results reported for the paper. The authors give some examples of the images generated from this setup. They note that these are selected to show the success of the model, not give a balanced representation of how good it is:<br />
<br />
[[File:Figure_1_AC-GAN.JPG|thumb|1000px|center|alt=(Odena et al., 2016) Figure 1: Selected images generated by the AC-GAN model for the ImageNet dataset.|(Odena et al., 2016) Figure 1: Selected images generated by the AC-GAN model for the ImageNet dataset.]]<br />
<br />
The apparent value of the model is its ability to generate high-resolution samples with global coherence. From the images given above one must acknowledge they are impressively realistic. The authors link to a [https://photos.google.com/share/AF1QipPzTToH3wxrKoF8l5nWvSNz7D_oS-KB3YMQ6Ji-4XK3AtgJmlb5QCdqRQLqLSjfkw?key=YnhJb2UyMnZwb01oZ2xhUjBraE9fdWU1VVpNZTVB full set] of images generated from every ImageNet class as well. As one would expect from their acknowledgement that images above are selected to be most impressive, not all samples at the linked page are quite so coherent.<br />
<br />
[[File:Figure_9_(Bottom).JPG|thumb|400px|right|alt=(Odena et al., 2016) Figure 9 (Bottom): Each column is a different class. Each row is generated by a different latent encoding $z$.|(Odena et al., 2016) Figure 9 (Bottom): Each column is a different class. Each row is generated by a different latent encoding $z$.]]<br />
<br />
The authors compare their model with state-of-the-art results from [[#References| Salimans et al. (2016)]] on the CIFAR-10 dataset at a 32 x 32 resolution. To score the two models they use Inception Score instead of log-likelihood, which they claim is too inaccurate to be reported. Their model achieves a score of $8.25 \pm 0.07$ versus the previous state-of-the-art of $8.09 \pm 0.07$.<br />
<br />
[[#References|Odena et al. (2016)]] argue that the class conditional generator allows $G$ to learn a representation of $Z$ independent of $C$ in section 3, and give evidence of the claim later in section 4.5 by showing that images generated with a fixed latent vector $z$ but different class labels $c$ have similar global structure (e.g. orientation of the subject) but the subjects (bird species) vary according to the label. Interestingly, the background (especially in the top row) also varies with the class label. This can possibly be attributed to the bird species coming from different areas, hence a seagull might be expected to have an ocean background. Clearly here the model benefits from the fact that the authors grouped similar classes together. A more interesting analysis might show the same comparison between different classes, such as birds and forklifts, to see how global structure is encoded across them.<br />
<br />
The authors also include a discussion of whether their model is overfitting the training data. Their first test is to find the nearest neighbour in the training set of a generated image by the L1 measure in pixel space and visually compare the two images. This is a fairly naive approach, since the L1 loss in pixel space is extremely unlikely to identify whether two images are perceptually similar. Here would have been a good place to use the MS-SSIM metric to identify the nearest neighbours, since it is intended to measure perceptual similarity. The images they provide from this analysis are below.<br />
<br />
[[File:Figure_8.JPG|thumb|500px|center|alt=(Odena et al., 2016) Figure 8: Images generated by the AC-GAN model (top) and their nearest neighbour (L1 measure in pixel space) in the ImageNet dataset (bottom).|(Odena et al., 2016) Figure 8: Images generated by the AC-GAN model (top) and their nearest neighbour (L1 measure in pixel space) in the ImageNet dataset (bottom).]]<br />
<br />
A second check they make is that interpolating between two generated images in the latent space does not result in any discrete transitions or holes in the image interpolation. Such results would be indicative of overfitting. The <br />
images they give as evidence that this is not the case are below.<br />
<br />
[[File:Figure_9_(Top).JPG|thumb|1000px|center|alt=(Odena et al., 2016) Figure 9 (Top): Latent space interpolations of the AC-GAN model.|(Odena et al., 2016) Figure 9 (Top): Latent space interpolations of the AC-GAN model.]]<br />
<br />
=== Image Diversity Metric ===<br />
<br />
Another result the authors report is the performance of their image diversity metric. It is difficult to evaluate quantitatively, but visually we see that the scores do appear to capture the perceptual diversity of the generated class. For example the 'artichoke' generator appears to have collapsed, and has a high score, while the 'promontory' generator seems fairly diverse and has a low score:<br />
<br />
[[File:Figure_3.JPG|thumb|700px|center|alt=(Odena et al., 2016) Figure 3: Mean MS-SSIM scores and sample images from selected generated (top) and real (bottom) ImageNet classes.|(Odena et al., 2016) Figure 3: Mean MS-SSIM scores and sample images from selected generated (top) and real (bottom) ImageNet classes.]]<br />
<br />
= Critique =<br />
=== Model ===<br />
<br />
A major attraction of this paper is the impressive quality of samples generated by the model. GANs often generate samples that are locally plausible but globally not realistic (e.g. a generated image of a dog has fur but the overall shape is not distinguishable). As we have seen in this critique, and as acknowledged by the authors, the most impressive samples are not representative of the model's overall performance.<br />
<br />
The model itself is not a very big advancement of the field. It combines two ideas that are both already prevalent in the research without any other justification than that it seems like a natural thing to do. As [https://openreview.net/forum?id=BkDDM04Ke other reviewers] have noted, investigating how much value the proposed model adds by comparing it with other models that only implement one (or neither) of the changes would have made this paper a slightly more interesting read.<br />
<br />
Another criticism I have about the paper is about how they report their results. To compare with [[#References|Salimans et al. (2016)]] they use Inception Score rather than log-likelihood, which they claim is the standard. Even if their model performed worse by that measure it ought to be included with the caveat they mentioned. The models are evaluated on a different dataset and at a lower spatial resolution than was used for the rest of the paper. By the Inception Score their results are better on average but might not be significantly different given how close they are. Finally, they did not apply the mean MS-SSIM score they developed in this paper to evaluate their model against [[#References|Salimans et al. (2016)]]. This would have been a natural point to make, but instead they generate four samples from each model as their evidence.<br />
<br />
An analysis the authors could have included that was touched upon but not explored in section 4.6 of the paper, and in the [[#Results|Results]] section of this summary, is how the similarity of the classes grouped in each sub-model impact the quality of generated samples. The example I gave above was to compare generated images with the same latent code but very different classes, such as birds and forklifts, to see how the global structure transferred across dissimilar classes.<br />
<br />
A last point to make about the model section is that the authors make some unsupported claims in their discussions of the model's properties. Specifically, they state that their modification to the standard GAN formulation appears to stabilize training but offer no evidence. Another example is their claim that "AC-GANs learn a representation for $z$ that is independent of class label". They site [[#References|Kingma et al (2014)]] as evidence of this. From my review of that paper it does not appear that the authors gave evidence for such a claim.<br />
<br />
=== GAN Quality Metrics ===<br />
<br />
The Inception accuracy metric proposed in this paper has the drawback that it is only applicable in a conditional GAN setting, since the standard GAN framework has no ground-truth labels. It is also true that using a pre-trained classifier is only a proxy for determining how much generated images look like the class they are meant to represent, since classifiers are not perfect. Consider the phenomenon of adversarial attacks on classifiers to see this point. However the advantages the authors list, that the Inception accuracy can be computed on a per-class basis and is easier to interpret than the Inception Score do have merit. The metric does make sense for the task the authors use it for.<br />
<br />
The same can be said for the mean MS-SSIM metric developed in this paper. Visually it appears to be a good indicator of diversity in the GAN's outputs. The authors claim that the mean MS-SSIM is a fast and easy-to-compute metric for perceptual variability and collapsing behaviour in a GAN. It is unclear how fast the metric can be computed since for each class the MS-SSIM has to be computed 100*99 times, once for each pair of images. The authors do not discuss how quickly it can actually be done.<br />
<br />
=== Experimental Results on GAN Properties ===<br />
<br />
The authors included three analyses which I have termed experiments. Of these, the first concluded that images generated at a higher resolution are more discriminable than images generated at lower resolutions, even when they have been resized to be comparable. This does not seem like a very revolutionary conclusion. For one thing, the space of lower resolution images in contained in the higher resolution space. In essence the high resolution model could generate lower resolution images by setting blocks of 4 pixels to the same intensity. It seems unsurprising then that the lower resolution is less discriminable on average. Another reason for this could be that the high resolution model has more parameters, and is trained on higher resolution data, so it has more information with which to reconstruct class information. Finally the authors give a graph of accuracies to show this property, and the average line appears compelling, however the standard errors about the lines suggest they may not be significantly different.<br />
<br />
The second experiment is on the interaction between the Inception accuracy and mean MS-SSIM metric. The author found that they are negatively correlated, and thus that classes that are high quality also tend to be diverse. This is contrary to prevailing wisdom, and since the correlation between them is weak, it appears that it may be only a fluke of the metrics.<br />
<br />
The final experiment is on the effect of class splits on image diversity. The authors found that increasing the number of classes handled by each model reduced the diversity of generated images. They make the claim at the beginning of the paper that they show the number of classes is what makes ImageNet synthesis difficult for GANs. This analysis does point in that direction but is not quite conclusive about the issue. Another analysis they could have included towards showing this is how their Inception accuracy metric and the Inception Score are affected by the number of class splits in their model.<br />
<br />
= Conclusion =<br />
<br />
This paper's main contributions were to introduce a slight variation on previous GAN models, as well as two metrics that can be used to assess the quality of generated images. The modified GAN, dubbed the Auxiliary Classifier GAN, was shown to produce high quality, high resolution samples from ImageNet, but not consistently The authors could have done more to show why their proposed architecture was an improvement over previous methods.<br />
<br />
The metrics introduced are both fairly straightforward and appear to function as they are intended. This being said, the authors could have used them more consistently throughout the paper (such as using the MS-SSIM to find nearest neighbours instead of the L1 pixel space loss). This paper was accepted to ICML 2017 but rejected by ICLR 2018 due to the incremental nature of the model development and the ad hoc nature of the other analyses presented above.<br />
<br />
= References =<br />
# Odena, A., Olah, C., & Shlens, J. (2016). Conditional image synthesis with auxiliary classifier gans. arXiv preprint [http://proceedings.mlr.press/v70/odena17a.html arXiv:1610.09585].<br />
# Goodfellow, I., Pouget-Abadie, J., Mirza, M., Xu, B., Warde-Farley, D., Ozair, S., ... & Bengio, Y. (2014). Generative adversarial nets. In Advances in neural information processing systems (pp. 2672-2680).<br />
# Salimans, T., Goodfellow, I., Zaremba, W., Cheung, V., Radford, A., & Chen, X. (2016). Improved techniques for training gans. In Advances in Neural Information Processing Systems (pp. 2234-2242).<br />
# Karras, T., Aila, T., Laine, S., & Lehtinen, J. (2017). Progressive Growing of GANs for Improved Quality, Stability, and Variation. arXiv preprint [https://arxiv.org/abs/1710.10196 arXiv:1710.10196].<br />
# Mirza, M., & Osindero, S. (2014). Conditional generative adversarial nets. arXiv preprint [https://arxiv.org/abs/1411.1784 arXiv:1411.1784].<br />
# van den Oord, A., Kalchbrenner, N., Espeholt, L., Vinyals, O., & Graves, A. (2016). Conditional image generation with pixelcnn decoders. In Advances in Neural Information Processing Systems (pp. 4790-4798).<br />
# Reed, S. E., Akata, Z., Mohan, S., Tenka, S., Schiele, B., & Lee, H. (2016). Learning what and where to draw. In Advances in Neural Information Processing Systems (pp. 217-225).<br />
# Szegedy, C., Liu, W., Jia, Y., Sermanet, P., Reed, S., Anguelov, D., ... & Rabinovich, A. (2015). Going deeper with convolutions. In Proceedings of the IEEE conference on computer vision and pattern recognition (pp. 1-9).<br />
# Sutskever, I., Vinyals, O., & Le, Q. V. (2014). Sequence to sequence learning with neural networks. In Advances in neural information processing systems (pp. 3104-3112).<br />
# Ramsundar, B., Kearnes, S., Riley, P., Webster, D., Konerding, D., & Pande, V. (2015). Massively multitask networks for drug discovery. arXiv preprint [https://arxiv.org/abs/1502.02072 arXiv:1502.02072]<br />
# Odena, A. (2016). Semi-supervised learning with generative adversarial networks. arXiv preprint [https://arxiv.org/abs/1606.01583 arXiv:1606.01583].<br />
# Theis, L., Oord, A. V. D., & Bethge, M. (2015). A note on the evaluation of generative models. arXiv preprint [https://arxiv.org/abs/1511.01844 arXiv:1511.01844].<br />
# Wang, Z., Simoncelli, E. P., & Bovik, A. C. (2003, November). Multiscale structural similarity for image quality assessment. In Signals, Systems and Computers, 2004. Conference Record of the Thirty-Seventh Asilomar Conference on (Vol. 2, pp. 1398-1402). IEEE.<br />
# Russakovsky, O., Deng, J., Su, H., Krause, J., Satheesh, S., Ma, S., ... & Berg, A. C. (2015). Imagenet large scale visual recognition challenge. International Journal of Computer Vision, 115(3), 211-252.</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Conditional_Image_Synthesis_with_Auxiliary_Classifier_GANs&diff=31290Conditional Image Synthesis with Auxiliary Classifier GANs2017-11-23T18:05:43Z<p>Nkadiyar: </p>
<hr />
<div>'''Abstract:''' "In this paper we introduce new methods for the improved training of generative adversarial networks (GANs) for image synthesis. We construct a variant of GANs employing label conditioning that results in 128×128 resolution image samples exhibiting global coherence. We expand on previous work for image quality assessment to provide two new analyses for assessing the discriminability and diversity of samples from class-conditional image synthesis models. These analyses demonstrate that high resolution samples provide class information not present in low resolution samples. Across 1000 ImageNet classes, 128×128 samples are more than twice as discriminable as artificially resized 32×32 samples. In addition, 84.7% of the classes have samples exhibiting diversity comparable to real ImageNet data." [[#References | (Odena et al., 2016)]]<br />
<br />
= Introduction =<br />
<br />
=== Motivation ===<br />
The authors introduce a GAN architecture for generating high resolution images from the ImageNet dataset. They show that this architecture makes it possible to split the generation process into many sub-models. They further suggest that GANs have trouble generating globally coherent images, and that this architecture is responsible for the coherence of their samples. They experimentally demonstrate that generating higher resolution images allow the model to encode more class-specific information, making them more visually discriminable than lower resolution images even after they have been resized to the same resolution.<br />
<br />
The second half of the paper introduces metrics for assessing visual discriminability and diversity of synthesized images. The discussion of image diversity in particular is important due to the tendency for GANs to 'collapse' to only produce one image that best fools the discriminator [[#References|(Goodfellow et al., 2014)]].<br />
<br />
=== Previous Work ===<br />
<br />
Of all image synthesis methods (e.g. variational autoencoders, autoregressive models, invertible density estimators), GANs have become one of the most popular and successful due to their flexibility and the ease with which they can be sampled from. A standard GAN framework pits a generative model $G$ against a discriminative adversary $D$. The goal of $G$ is to learn a mapping from a latent space $Z$ to a real space $X$ to produce examples (generally images) indistinguishable from training data. The goal of the $D$ is to iteratively learn to predict when a given input image is from the training set or a synthesized image from $G$. Jointly the models are trained to solve the game-theoretical minimax problem, as defined by [[#References|Goodfellow et al. (2014)]]: $$\underset{G}{\text{min }}\underset{D}{\text{max }}V(G,D)=\mathbb{E}_{X\sim p_{data}(x)}[log(D(X))]+\mathbb{E}_{Z\sim p_{Z}(z)}[log(1-D(G(Z)))]$$<br />
<br />
While this initial framework has clearly demonstrated great potential, other authors have proposed changes to the method to improve it. Many such papers propose changes to the training process [[#References|(Salimans et al., 2016)]][[#References|(Karras et al., 2017)]], which is notoriously difficult for some problems. Others propose changes to the model itself. [[#References|Mirza & Osindero (2014)]] augment the model by supplying the class of observations to both the generator and discriminator to produce class-conditional samples. According to [[STAT946F17/Conditional Image Generation with PixelCNN Decoders|van den Oord et al. (2016)]], conditioning image generation on classes can greatly improve their quality. Other authors have explored using even richer side information in the generation process with good results [[Learning What and Where to Draw|(Reed et al., 2016)]].<br />
<br />
Another model modification relevant to this paper is to force the discriminator network to reconstruct side information by adding an auxiliary network to classify generated (and real) images. The authors make the claim that forcing a model to perform additional tasks is known to improve performance on the original task [[#References|(Szegedy et al., 2014)]][[#References|(Sutskever et al., 2014)]][[#References|(Ramsundar et al., 2016)]]. They further suggest that using pre-trained image classifiers (rather than classifiers trained on both real and generated images) could improve results over and above what is shown in this paper.<br />
<br />
= Contributions =<br />
<br />
The contributions of this paper are in three main areas. First, the authors propose slight changes to previously existing GAN architectures, resulting in a model capable of generating samples of impressive quality. Second, the authors propose two metrics to assess the quality of samples generated from a GAN. Lastly, they present empirical results on GANs which are of some interest.<br />
<br />
== Model ==<br />
<br />
The authors propose an auxiliary classifier GAN (AC-GAN) which is a slight variation on previous architectures. Like [[#References|Mirza & Osindero (2014)]], the generator takes the image class to be generated as input in addition to the latent encoding $Z$. Like [[#References|Odena (2016)]] and [[#References|Salimans et al. (2016)]], the discriminator is trained to predict not only whether an observation is real or fake, but to classify each observation as well. The marginal contribution of this paper is to combine these in one model.<br />
<br />
Formally, let $C\sim p_c$ represent the target class label of each generated observation and $Z$ represent the usual noise vector from the latent space. Then the generator function takes both as argument to produce image samples: $X_{fake}=G(c,z)$.The discriminator gives a probability distribution over the source $S$ (real or fake) of the image as well as the class label $C$ being generated. $$D(X)=<P(S|X),P(C|X)>$$<br />
<br />
The objective function for the model thus has two parts, one corresponding to the source $L_S$ and the other to the class $L_C$. $D$ is trained to maximize $L_S + L_C$, while $G$ is trained to maximize $L_C-L_S$. Using the notation of Goodfellow et al. (2014), the loss terms are defined as:<br />
$$L_S=\mathbb{E}_{X\sim p_{data}(x)}[log(P(S=real|X))]+\mathbb{E}_{C,Z\sim p_{C,Z}(c,z)}[log(P(S=fake|G(C,Z)))]$$<br />
$$L_C=\mathbb{E}_{X\sim p_{data}(x)}[log(P(C|X))]+\mathbb{E}_{C,Z\sim p_{C,Z}(c,z)}[log(P(C|G(C,Z)))]$$<br />
<br />
Because G accepts both $C$ and $Z$ as arguments, it is able to learn a mapping $Z\rightarrow X$ that is independent of $C$. The authors argue that all class-specific information should be represented by $C$, allowing $Z$ to represent other factors such as pose, background, etc.<br />
<br />
Lastly the authors split the generation process into many class-specific submodels. They point out that the structure of their model permits this split, though it should technically be possible for even the standard GAN framework by dividing the training data into groups according to their known class labels.<br />
<br />
The changes above result in a model capable of generating (some) image samples with both high resolution and global coherence.<br />
<br />
== GAN Quality Metrics ==<br />
<br />
A much larger part of the paper is spent on measuring the quality of a GAN's output. As the authors say, evaluating a generative model's quality is difficult due to a large number of probabilistic measures (such as average log-likelihood, Parzen window estimates, and visual fidelity [[#References| (Theis et al., 2015)]]) and "a lack of a perceptually meaningful image similarity metric".<br />
<br />
=== Image Discriminability Metric ===<br />
The authors develop two metrics in this paper to address these shortcomings. The first of these is a discriminability metric, the goal of which is to assess the degree to which generated images are identifiable as the class they are meant to represent. Ideally a team of non-expert humans could handle this, but the difficulty of such an approach makes the need for an automated metric apparent. The metric proposed by the authors is to measure the accuracy of a pre-trained image classifier trained on the pristine training data. For this they select a [https://github.com/openai/improved-gan/ modified version of Inception-v3][[#References|(Szegedy et al., 2015)]].<br />
<br />
Other metrics already exist for assessing image quality, the most popular of which is probably the Inception Score [[#References| (Salimans et al., 2016)]]. The authors list two main advantages of their approach over the Inception Score. The first is that accuracy figures are easier to interpret than the Inception Score, which is fairly self-evident. The second advantage of using Inception accuracy instead of Inception Score is that Inception accuracy may be calculated for individual classes, giving a better picture of where the model is strong and where it is weak.<br />
<br />
=== Image Diversity Metric ===<br />
<br />
The second metric proposed by the authors measures the diversity of the generated images. As mentioned above, image diversity is an important quality in a GAN, a common failure they suffer from is 'collapsing', where the generator learns to only output one image that is good at fooling the discriminator [[#References|(Goodfellow et al., 2014)]][[#References|(Salimans et al., 2016)]]. The metric proposed in this section is intended to be complementary to the Inception accuracy, as Inception accuracy would not detect generator collapse.<br />
<br />
For their diversity metric, the authors co-opt an existing metric used to measure the similarity between two images: multi-scale structural similarity (MS-SSIM)[[#References|(Wang et al., 2003)]]. The authors do not go into detail about how the MS-SSIM measure is calculated, except to say that it is one of the more successful ways to predict human perceptual similarity judgements. It can take values on the interval $[0,1]$, and higher values indicate the two images being compared a perceptually more similar. For images generated from a GAN then, the metric should ideally be low, as diversity is desired.<br />
<br />
The authors contribution is to use this metric to assess diversity of a GAN's output. It is a pairwise comparison between two images, so their solution is to compare 100 images (that is $100\cdot 99$ paired comparisons) from each class and take the mean MS-SSIM score.<br />
<br />
The authors make two points about their use of this metric. First, the way they apply the metric is different from how it was originally intended to be used. It is possible that it will not behave as desired because of this. As evidence to the contrary, they state that:<br />
# Visually the metric seems to work. Pairs with high MS-SSIM seem more similar.<br />
# Comparisons are only made between images in the same class, keeping their application of the metric closer to its original use case of measuring the quality of compression algorithms.<br />
# The metric is not saturated. Scores on their generated data vary across the unit interval. If scores were all very close to zero the metric would not be much use.<br />
<br />
The second point they raise is that the mean MS-SSIM metric is not intended as a proxy for entropy of the generator distribution in pixel space. That measure is hard to compute, and in any case is sensitive to trivial changes in the pixels, whereas the true intention of this metric is to measure perceptual similarity.<br />
<br />
== Experimental Results on GAN Properties ==<br />
<br />
[[File:Figure_2_(Bottom).JPG|thumb|500px|right|alt=(Odena et al., 2016) Figure 2: (Left) Inception accuracy (y-axis) of two generators with resolution 128 x 128 (red) and 64 x 64 (blue). Images are resized to the same spatial resolution (x-axis). (Right) Class accuracies from the 128 x 128 AC-GAN at full resolution (x-axis) and downsized to 32 x 32 (y-axis).|(Odena et al., 2016) Figure 2: (Left) Inception accuracy (y-axis) of two generators with resolution 128 x 128 (red) and 64 x 64 (blue). Images are resized to the same spatial resolution (x-axis). (Right) Class accuracies from the 128 x 128 AC-GAN at full resolution (x-axis) and downsized to 32 x 32 (y-axis).]]<br />
<br />
The authors conduct several experiments on their model and proposed metrics. These are summarized in this section.<br />
<br />
=== Higher Resolution Images are more Discriminable ===<br />
<br />
As it is one of the main attractions of this paper, the authors investigate how generating samples at different resolutions affects their discriminability. To achieve this two models are trained, one that generates 64 x 64 resolution images and one that generates 128 x 128 resolution images. These images can be rescaled using bilinear interpolation to make them directly comparable. The authors find that the 128 x 128 AC-GAN (on average) achieves higher discriminability (per the Inception accuracy metric [[#Image Discriminability Metric|introduced above]]) at all resized resolutions. The authors claim that theirs is the first work to investigate how much an image generator is 'making use of its given output resolution'.<br />
<br />
=== Generated Images are both Diverse and Discriminable ===<br />
<br />
A second experiment conducted in the paper aims to investigate how the two metrics they propose interact with each other. They simply calculate the Inception accuracy and mean MS-SSIM score for a batch of generated images from every class in their data and report the correlation between the scores. They find the scores are anti-correlated with $\rho=-0.16$. Because the mean MS-SSIM metric is low for diverse samples, they conclude that accuracy and diversity are actually positively correlated, which is contradictory to the hypothesis that GANs that collapse achieve better sample quality.<br />
<br />
[[File:Figure_10.JPG|thumb|300px|right|alt=(Odena et al., 2016) Figure 10: Mean MS-SSIM scores for 10 ImageNet classes (y-axis) plotted against the number of classes handled by each sub-model (x-axis).|(Odena et al., 2016) Figure 10: Mean MS-SSIM scores for 10 ImageNet classes (y-axis) plotted against the number of classes handled by each sub-model (x-axis).]]<br />
<br />
=== Effect of Class Splits on Image Sample Quality ===<br />
<br />
The final experiment conducted in the paper investigates how the number of classes a fixed GAN model has to generate impacts the diversity of the generated images. In their main model the authors split the data such that each AC-GAN only had to generate images from 10 classes. Here they experimented with changing the number of classes each sub-model has to generate while holding the architecture of the sub-models fixed. They report the mean MS-SSIM score of the generated images from the original classes (when the model had only 10 to generate) at each split level. Perhaps unsurprisingly, the diversity of the outputs drops as the number of classes the model has to handle increases. Giving the model more parameters to handle the larger (and more diverse) set of classes might possibly eliminate this effect.<br />
<br />
Finally, they state that they were unable to get a regular GAN to converge on the task of generating images from one class per model. This could be due to the difficulty of training GANs and the limited amount of training data available for each class rather than any theoretical property of class splits.<br />
<br />
= Results =<br />
<br />
=== Model ===<br />
The authors apply their AC-GAN model to the ImageNet [[#References|(Russakovsky et al., 2015)]] dataset. The data have 1000 classes which the authors split into 100 groups of 10. An AC-GAN model is trained on each group of 10 to give results reported for the paper. The authors give some examples of the images generated from this setup. They note that these are selected to show the success of the model, not give a balanced representation of how good it is:<br />
<br />
[[File:Figure_1_AC-GAN.JPG|thumb|1000px|center|alt=(Odena et al., 2016) Figure 1: Selected images generated by the AC-GAN model for the ImageNet dataset.|(Odena et al., 2016) Figure 1: Selected images generated by the AC-GAN model for the ImageNet dataset.]]<br />
<br />
The apparent value of the model is its ability to generate high-resolution samples with global coherence. From the images given above one must acknowledge they are impressively realistic. The authors link to a [https://photos.google.com/share/AF1QipPzTToH3wxrKoF8l5nWvSNz7D_oS-KB3YMQ6Ji-4XK3AtgJmlb5QCdqRQLqLSjfkw?key=YnhJb2UyMnZwb01oZ2xhUjBraE9fdWU1VVpNZTVB full set] of images generated from every ImageNet class as well. As one would expect from their acknowledgement that images above are selected to be most impressive, not all samples at the linked page are quite so coherent.<br />
<br />
[[File:Figure_9_(Bottom).JPG|thumb|400px|right|alt=(Odena et al., 2016) Figure 9 (Bottom): Each column is a different class. Each row is generated by a different latent encoding $z$.|(Odena et al., 2016) Figure 9 (Bottom): Each column is a different class. Each row is generated by a different latent encoding $z$.]]<br />
<br />
The authors compare their model with state-of-the-art results from [[#References| Salimans et al. (2016)]] on the CIFAR-10 dataset at a 32 x 32 resolution. To score the two models they use Inception Score instead of log-likelihood, which they claim is too inaccurate to be reported. Their model achieves a score of $8.25 \pm 0.07$ versus the previous state-of-the-art of $8.09 \pm 0.07$.<br />
<br />
[[#References|Odena et al. (2016)]] argue that the class conditional generator allows $G$ to learn a representation of $Z$ independent of $C$ in section 3, and give evidence of the claim later in section 4.5 by showing that images generated with a fixed latent vector $z$ but different class labels $c$ have similar global structure (e.g. orientation of the subject) but the subjects (bird species) vary according to the label. Interestingly, the background (especially in the top row) also varies with the class label. This can possibly be attributed to the bird species coming from different areas, hence a seagull might be expected to have an ocean background. Clearly here the model benefits from the fact that the authors grouped similar classes together. A more interesting analysis might show the same comparison between different classes, such as birds and forklifts, to see how global structure is encoded across them.<br />
<br />
The authors also include a discussion of whether their model is overfitting the training data. Their first test is to find the nearest neighbour in the training set of a generated image by the L1 measure in pixel space and visually compare the two images. This is a fairly naive approach, since the L1 loss in pixel space is extremely unlikely to identify whether two images are perceptually similar. Here would have been a good place to use the MS-SSIM metric to identify the nearest neighbours, since it is intended to measure perceptual similarity. The images they provide from this analysis are below.<br />
<br />
[[File:Figure_8.JPG|thumb|500px|center|alt=(Odena et al., 2016) Figure 8: Images generated by the AC-GAN model (top) and their nearest neighbour (L1 measure in pixel space) in the ImageNet dataset (bottom).|(Odena et al., 2016) Figure 8: Images generated by the AC-GAN model (top) and their nearest neighbour (L1 measure in pixel space) in the ImageNet dataset (bottom).]]<br />
<br />
A second check they make is that interpolating between two generated images in the latent space does not result in any discrete transitions or holes in the image interpolation. Such results would be indicative of overfitting. The <br />
images they give as evidence that this is not the case are below.<br />
<br />
[[File:Figure_9_(Top).JPG|thumb|1000px|center|alt=(Odena et al., 2016) Figure 9 (Top): Latent space interpolations of the AC-GAN model.|(Odena et al., 2016) Figure 9 (Top): Latent space interpolations of the AC-GAN model.]]<br />
<br />
=== Image Diversity Metric ===<br />
<br />
Another result the authors report is the performance of their image diversity metric. It is difficult to evaluate quantitatively, but visually we see that the scores do appear to capture the perceptual diversity of the generated class. For example the 'artichoke' generator appears to have collapsed, and has a high score, while the 'promontory' generator seems fairly diverse and has a low score:<br />
<br />
[[File:Figure_3.JPG|thumb|700px|center|alt=(Odena et al., 2016) Figure 3: Mean MS-SSIM scores and sample images from selected generated (top) and real (bottom) ImageNet classes.|(Odena et al., 2016) Figure 3: Mean MS-SSIM scores and sample images from selected generated (top) and real (bottom) ImageNet classes.]]<br />
<br />
= Critique =<br />
=== Model ===<br />
<br />
A major attraction of this paper is the impressive quality of samples generated by the model. GANs often generate samples that are locally plausible but globally not realistic (e.g. a generated image of a dog has fur but the overall shape is not distinguishable). As we have seen in this critique, and as acknowledged by the authors, the most impressive samples are not representative of the model's overall performance.<br />
<br />
The model itself is not a very big advancement of the field. It combines two ideas that are both already prevalent in the research without any other justification than that it seems like a natural thing to do. As [https://openreview.net/forum?id=BkDDM04Ke other reviewers] have noted, investigating how much value the proposed model adds by comparing it with other models that only implement one (or neither) of the changes would have made this paper a slightly more interesting read.<br />
<br />
Another criticism I have about the paper is about how they report their results. To compare with [[#References|Salimans et al. (2016)]] they use Inception Score rather than log-likelihood, which they claim is the standard. Even if their model performed worse by that measure it ought to be included with the caveat they mentioned. The models are evaluated on a different dataset and at a lower spatial resolution than was used for the rest of the paper. By the Inception Score their results are better on average but might not be significantly different given how close they are. Finally, they did not apply the mean MS-SSIM score they developed in this paper to evaluate their model against [[#References|Salimans et al. (2016)]]. This would have been a natural point to make, but instead they generate four samples from each model as their evidence.<br />
<br />
An analysis the authors could have included that was touched upon but not explored in section 4.6 of the paper, and in the [[#Results|Results]] section of this summary, is how the similarity of the classes grouped in each sub-model impact the quality of generated samples. The example I gave above was to compare generated images with the same latent code but very different classes, such as birds and forklifts, to see how the global structure transferred across dissimilar classes.<br />
<br />
A last point to make about the model section is that the authors make some unsupported claims in their discussions of the model's properties. Specifically, they state that their modification to the standard GAN formulation appears to stabilize training but offer no evidence. Another example is their claim that "AC-GANs learn a representation for $z$ that is independent of class label". They site [[#References|Kingma et al (2014)]] as evidence of this. From my review of that paper it does not appear that the authors gave evidence for such a claim.<br />
<br />
=== GAN Quality Metrics ===<br />
<br />
The Inception accuracy metric proposed in this paper has the drawback that it is only applicable in a conditional GAN setting, since the standard GAN framework has no ground-truth labels. It is also true that using a pre-trained classifier is only a proxy for determining how much generated images look like the class they are meant to represent, since classifiers are not perfect. Consider the phenomenon of adversarial attacks on classifiers to see this point. However the advantages the authors list, that the Inception accuracy can be computed on a per-class basis and is easier to interpret than the Inception Score do have merit. The metric does make sense for the task the authors use it for.<br />
<br />
The same can be said for the mean MS-SSIM metric developed in this paper. Visually it appears to be a good indicator of diversity in the GAN's outputs. The authors claim that the mean MS-SSIM is a fast and easy-to-compute metric for perceptual variability and collapsing behaviour in a GAN. It is unclear how fast the metric can be computed since for each class the MS-SSIM has to be computed 100*99 times, once for each pair of images. The authors do not discuss how quickly it can actually be done.<br />
<br />
=== Experimental Results on GAN Properties ===<br />
<br />
The authors included three analyses which I have termed experiments. Of these, the first concluded that images generated at a higher resolution are more dscriminable than images generated at lower resolutions, even when they have been resized to be comparable. This does not seem like a very revolutionary conclusion. For one thing, the space of lower resolution images in contained in the higher resolution space. In essence the high resolution model could generate lower resolution images by setting blocks of 4 pixels to the same intensity. It seems unsurprising then that the lower resolution is less discriminable on average. Another reason for this could be that the high resolution model has more parameters, and is trained on higher resolution data, so it has more information with which to reconstruct class information. Finally the authors give a graph of accuracies to show this property, and the average line appears compelling, however the standard errors about the lines suggest they may not be significantly different.<br />
<br />
The second experiment is on the interaction between the Inception accuracy and mean MS-SSIM metric. The author found that they are negatively correlated, and thus that classes that are high quality also tend to be diverse. This is contrary to prevailing wisdom, and since the correlation between them is weak, it appears that it may be only a fluke of the metrics.<br />
<br />
The final experiment is on the effect of class splits on image diversity. The authors found that increasing the number of classes handled by each model reduced the diversity of generated images. They make the claim at the beginning of the paper that they show the number of classes is what makes ImageNet synthesis difficult for GANs. This analysis does point in that direction but is not quite conclusive about the issue. Another analysis they could have included towards showing this is how their Inception accuracy metric and the Inception Score are affected by the number of class splits in their model.<br />
<br />
= Conclusion =<br />
<br />
This paper's main contributions were to introduce a slight variation on previous GAN models, as well as two metrics that can be used to assess the quality of generated images. The modified GAN, dubbed the Auxiliary Classifier GAN, was shown to produce high quality, high resolution samples from ImageNet, but not consistently The authors could have done more to show why their proposed architecture was an improvement over previous methods.<br />
<br />
The metrics introduced are both fairly straightforward and appear to function as they are intended. This being said, the authors could have used them more consistently throughout the paper (such as using the MS-SSIM to find nearest neighbours instead of the L1 pixel space loss). This paper was accepted to ICML 2017 but rejected by ICLR 2018 due to the incremental nature of the model development and the ad hoc nature of the other analyses presented above.<br />
<br />
= References =<br />
# Odena, A., Olah, C., & Shlens, J. (2016). Conditional image synthesis with auxiliary classifier gans. arXiv preprint [http://proceedings.mlr.press/v70/odena17a.html arXiv:1610.09585].<br />
# Goodfellow, I., Pouget-Abadie, J., Mirza, M., Xu, B., Warde-Farley, D., Ozair, S., ... & Bengio, Y. (2014). Generative adversarial nets. In Advances in neural information processing systems (pp. 2672-2680).<br />
# Salimans, T., Goodfellow, I., Zaremba, W., Cheung, V., Radford, A., & Chen, X. (2016). Improved techniques for training gans. In Advances in Neural Information Processing Systems (pp. 2234-2242).<br />
# Karras, T., Aila, T., Laine, S., & Lehtinen, J. (2017). Progressive Growing of GANs for Improved Quality, Stability, and Variation. arXiv preprint [https://arxiv.org/abs/1710.10196 arXiv:1710.10196].<br />
# Mirza, M., & Osindero, S. (2014). Conditional generative adversarial nets. arXiv preprint [https://arxiv.org/abs/1411.1784 arXiv:1411.1784].<br />
# van den Oord, A., Kalchbrenner, N., Espeholt, L., Vinyals, O., & Graves, A. (2016). Conditional image generation with pixelcnn decoders. In Advances in Neural Information Processing Systems (pp. 4790-4798).<br />
# Reed, S. E., Akata, Z., Mohan, S., Tenka, S., Schiele, B., & Lee, H. (2016). Learning what and where to draw. In Advances in Neural Information Processing Systems (pp. 217-225).<br />
# Szegedy, C., Liu, W., Jia, Y., Sermanet, P., Reed, S., Anguelov, D., ... & Rabinovich, A. (2015). Going deeper with convolutions. In Proceedings of the IEEE conference on computer vision and pattern recognition (pp. 1-9).<br />
# Sutskever, I., Vinyals, O., & Le, Q. V. (2014). Sequence to sequence learning with neural networks. In Advances in neural information processing systems (pp. 3104-3112).<br />
# Ramsundar, B., Kearnes, S., Riley, P., Webster, D., Konerding, D., & Pande, V. (2015). Massively multitask networks for drug discovery. arXiv preprint [https://arxiv.org/abs/1502.02072 arXiv:1502.02072]<br />
# Odena, A. (2016). Semi-supervised learning with generative adversarial networks. arXiv preprint [https://arxiv.org/abs/1606.01583 arXiv:1606.01583].<br />
# Theis, L., Oord, A. V. D., & Bethge, M. (2015). A note on the evaluation of generative models. arXiv preprint [https://arxiv.org/abs/1511.01844 arXiv:1511.01844].<br />
# Wang, Z., Simoncelli, E. P., & Bovik, A. C. (2003, November). Multiscale structural similarity for image quality assessment. In Signals, Systems and Computers, 2004. Conference Record of the Thirty-Seventh Asilomar Conference on (Vol. 2, pp. 1398-1402). IEEE.<br />
# Russakovsky, O., Deng, J., Su, H., Krause, J., Satheesh, S., Ma, S., ... & Berg, A. C. (2015). Imagenet large scale visual recognition challenge. International Journal of Computer Vision, 115(3), 211-252.</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Conditional_Image_Synthesis_with_Auxiliary_Classifier_GANs&diff=31289Conditional Image Synthesis with Auxiliary Classifier GANs2017-11-23T18:04:58Z<p>Nkadiyar: </p>
<hr />
<div>'''Abstract:''' "In this paper we introduce new methods for the improved training of generative adversarial networks (GANs) for image synthesis. We construct a variant of GANs employing label conditioning that results in 128×128 resolution image samples exhibiting global coherence. We expand on previous work for image quality assessment to provide two new analyses for assessing the discriminability and diversity of samples from class-conditional image synthesis models. These analyses demonstrate that high resolution samples provide class information not present in low resolution samples. Across 1000 ImageNet classes, 128×128 samples are more than twice as discriminable as artificially resized 32×32 samples. In addition, 84.7% of the classes have samples exhibiting diversity comparable to real ImageNet data." [[#References | (Odena et al., 2016)]]<br />
<br />
= Introduction =<br />
<br />
=== Motivation ===<br />
The authors introduce a GAN architecture for generating high resolution images from the ImageNet dataset. They show that this architecture makes it possible to split the generation process into many sub-models. They further suggest that GANs have trouble generating globally coherent images, and that this architecture is responsible for the coherence of their samples. They experimentally demonstrate that generating higher resolution images allow the model to encode more class-specific information, making them more visually discriminable than lower resolution images even after they have been resized to the same resolution.<br />
<br />
The second half of the paper introduces metrics for assessing visual discriminability and diversity of synthesized images. The discussion of image diversity in particular is important due to the tendency for GANs to 'collapse' to only produce one image that best fools the discriminator [[#References|(Goodfellow et al., 2014)]].<br />
<br />
=== Previous Work ===<br />
<br />
Of all image synthesis methods (e.g. variational autoencoders, autoregressive models, invertible density estimators), GANs have become one of the most popular and successful due to their flexibility and the ease with which they can be sampled from. A standard GAN framework pits a generative model $G$ against a discriminative adversary $D$. The goal of $G$ is to learn a mapping from a latent space $Z$ to a real space $X$ to produce examples (generally images) indistinguishable from training data. The goal of the $D$ is to iteratively learn to predict when a given input image is from the training set or a synthesized image from $G$. Jointly the models are trained to solve the game-theoretical minimax problem, as defined by [[#References|Goodfellow et al. (2014)]]: $$\underset{G}{\text{min }}\underset{D}{\text{max }}V(G,D)=\mathbb{E}_{X\sim p_{data}(x)}[log(D(X))]+\mathbb{E}_{Z\sim p_{Z}(z)}[log(1-D(G(Z)))]$$<br />
<br />
While this initial framework has clearly demonstrated great potential, other authors have proposed changes to the method to improve it. Many such papers propose changes to the training process [[#References|(Salimans et al., 2016)]][[#References|(Karras et al., 2017)]], which is notoriously difficult for some problems. Others propose changes to the model itself. [[#References|Mirza & Osindero (2014)]] augment the model by supplying the class of observations to both the generator and discriminator to produce class-conditional samples. According to [[STAT946F17/Conditional Image Generation with PixelCNN Decoders|van den Oord et al. (2016)]], conditioning image generation on classes can greatly improve their quality. Other authors have explored using even richer side information in the generation process with good results [[Learning What and Where to Draw|(Reed et al., 2016)]].<br />
<br />
Another model modification relevant to this paper is to force the discriminator network to reconstruct side information by adding an auxiliary network to classify generated (and real) images. The authors make the claim that forcing a model to perform additional tasks is known to improve performance on the original task [[#References|(Szegedy et al., 2014)]][[#References|(Sutskever et al., 2014)]][[#References|(Ramsundar et al., 2016)]]. They further suggest that using pre-trained image classifiers (rather than classifiers trained on both real and generated images) could improve results over and above what is shown in this paper.<br />
<br />
= Contributions =<br />
<br />
The contributions of this paper are in three main areas. First, the authors propose slight changes to previously existing GAN architectures, resulting in a model capable of generating samples of impressive quality. Second, the authors propose two metrics to assess the quality of samples generated from a GAN. Lastly, they present empirical results on GANs which are of some interest.<br />
<br />
== Model ==<br />
<br />
The authors propose an auxiliary classifier GAN (AC-GAN) which is a slight variation on previous architectures. Like [[#References|Mirza & Osindero (2014)]], the generator takes the image class to be generated as input in addition to the latent encoding $Z$. Like [[#References|Odena (2016)]] and [[#References|Salimans et al. (2016)]], the discriminator is trained to predict not only whether an observation is real or fake, but to classify each observation as well. The marginal contribution of this paper is to combine these in one model.<br />
<br />
Formally, let $C\sim p_c$ represent the target class label of each generated observation and $Z$ represent the usual noise vector from the latent space. Then the generator function takes both as argument to produce image samples: $X_{fake}=G(c,z)$.The discriminator gives a probability distribution over the source $S$ (real or fake) of the image as well as the class label $C$ being generated. $$D(X)=<P(S|X),P(C|X)>$$<br />
<br />
The objective function for the model thus has two parts, one corresponding to the source $L_S$ and the other to the class $L_C$. $D$ is trained to maximize $L_S + L_C$, while $G$ is trained to maximize $L_C-L_S$. Using the notation of Goodfellow et al. (2014), the loss terms are defined as:<br />
$$L_S=\mathbb{E}_{X\sim p_{data}(x)}[log(P(S=real|X))]+\mathbb{E}_{C,Z\sim p_{C,Z}(c,z)}[log(P(S=fake|G(C,Z)))]$$<br />
$$L_C=\mathbb{E}_{X\sim p_{data}(x)}[log(P(C|X))]+\mathbb{E}_{C,Z\sim p_{C,Z}(c,z)}[log(P(C|G(C,Z)))]$$<br />
<br />
Because G accepts both $C$ and $Z$ as arguments, it is able to learn a mapping $Z\rightarrow X$ that is independent of $C$. The authors argue that all class-specific information should be represented by $C$, allowing $Z$ to represent other factors such as pose, background, etc.<br />
<br />
Lastly the authors split the generation process into many class-specific submodels. They point out that the structure of their model permits this split, though it should technically be possible for even the standard GAN framework by dividing the training data into groups according to their known class labels.<br />
<br />
The changes above result in a model capable of generating (some) image samples with both high resolution and global coherence.<br />
<br />
== GAN Quality Metrics ==<br />
<br />
A much larger part of the paper is spent on measuring the quality of a GAN's output. As the authors say, evaluating a generative model's quality is difficult due to a large number of probabilistic measures (such as average log-likelihood, Parzen window estimates, and visual fidelity [[#References| (Theis et al., 2015)]]) and "a lack of a perceptually meaningful image similarity metric".<br />
<br />
=== Image Discriminability Metric ===<br />
The authors develop two metrics in this paper to address these shortcomings. The first of these is a discriminability metric, the goal of which is to assess the degree to which generated images are identifiable as the class they are meant to represent. Ideally a team of non-expert humans could handle this, but the difficulty of such an approach makes the need for an automated metric apparent. The metric proposed by the authors is to measure the accuracy of a pre-trained image classifier trained on the pristine training data. For this they select a [https://github.com/openai/improved-gan/ modified version of Inception-v3][[#References|(Szegedy et al., 2015)]].<br />
<br />
Other metrics already exist for assessing image quality, the most popular of which is probably the Inception Score [[#References| (Salimans et al., 2016)]]. The authors list two main advantages of their approach over the Inception Score. The first is that accuracy figures are easier to interpret than the Inception Score, which is fairly self-evident. The second advantage of using Inception accuracy instead of Inception Score is that Inception accuracy may be calculated for individual classes, giving a better picture of where the model is strong and where it is weak.<br />
<br />
=== Image Diversity Metric ===<br />
<br />
The second metric proposed by the authors measures the diversity of the generated images. As mentioned above, image diversity is an important quality in a GAN, a common failure they suffer from is 'collapsing', where the generator learns to only output one image that is good at fooling the discriminator [[#References|(Goodfellow et al., 2014)]][[#References|(Salimans et al., 2016)]]. The metric proposed in this section is intended to be complementary to the Inception accuracy, as Inception accuracy would not detect generator collapse.<br />
<br />
For their diversity metric, the authors co-opt an existing metric used to measure the similarity between two images: multi-scale structural similarity (MS-SSIM)[[#References|(Wang et al., 2003)]]. The authors do not go into detail about how the MS-SSIM measure is calculated, except to say that it is one of the more successful ways to predict human perceptual similarity judgements. It can take values on the interval $[0,1]$, and higher values indicate the two images being compared a perceptually more similar. For images generated from a GAN then, the metric should ideally be low, as diversity is desired.<br />
<br />
The authors contribution is to use this metric to assess diversity of a GAN's output. It is a pairwise comparison between two images, so their solution is to compare 100 images (that is $100\cdot 99$ paired comparisons) from each class and take the mean MS-SSIM score.<br />
<br />
The authors make two points about their use of this metric. First, the way they apply the metric is different from how it was originally intended to be used. It is possible that it will not behave as desired because of this. As evidence to the contrary, they state that:<br />
# Visually the metric seems to work. Pairs with high MS-SSIM seem more similar.<br />
# Comparisons are only made between images in the same class, keeping their application of the metric closer to its original use case of measuring the quality of compression algorithms.<br />
# The metric is not saturated. Scores on their generated data vary across the unit interval. If scores were all very close to zero the metric would not be much use.<br />
<br />
The second point they raise is that the mean MS-SSIM metric is not intended as a proxy for entropy of the generator distribution in pixel space. That measure is hard to compute, and in any case is sensitive to trivial changes in the pixels, whereas the true intention of this metric is to measure perceptual similarity.<br />
<br />
== Experimental Results on GAN Properties ==<br />
<br />
[[File:Figure_2_(Bottom).JPG|thumb|500px|right|alt=(Odena et al., 2016) Figure 2: (Left) Inception accuracy (y-axis) of two generators with resolution 128 x 128 (red) and 64 x 64 (blue). Images are resized to the same spatial resolution (x-axis). (Right) Class accuracies from the 128 x 128 AC-GAN at full resolution (x-axis) and downsized to 32 x 32 (y-axis).|(Odena et al., 2016) Figure 2: (Left) Inception accuracy (y-axis) of two generators with resolution 128 x 128 (red) and 64 x 64 (blue). Images are resized to the same spatial resolution (x-axis). (Right) Class accuracies from the 128 x 128 AC-GAN at full resolution (x-axis) and downsized to 32 x 32 (y-axis).]]<br />
<br />
The authors conduct several experiments on their model and proposed metrics. These are summarized in this section.<br />
<br />
=== Higher Resolution Images are more Discriminable ===<br />
<br />
As it is one of the main attractions of this paper, the authors investigate how generating samples at different resolutions affects their discriminability. To achieve this two models are trained, one that generates 64 x 64 resolution images and one that generates 128 x 128 resolution images. These images can be rescaled using bilinear interpolation to make them directly comparable. The authors find that the 128 x 128 AC-GAN (on average) achieves higher discriminability (per the Inception accuracy metric [[#Image Discriminability Metric|introduced above]]) at all resized resolutions. The authors claim that theirs is the first work to investigate how much an image generator is 'making use of its given output resolution'.<br />
<br />
=== Generated Images are both Diverse and Discriminable ===<br />
<br />
A second experiment conducted in the paper aims to investigate how the two metrics they propose interact with each other. They simply calculate the Inception accuracy and mean MS-SSIM score for a batch of generated images from every class in their data and report the correlation between the scores. They find the scores are anti-correlated with $\rho=-0.16$. Because the mean MS-SSIM metric is low for diverse samples, they conclude that accuracy and diversity are actually positively correlated, which is contradictory to the hypothesis that GANs that collapse achieve better sample quality.<br />
<br />
[[File:Figure_10.JPG|thumb|300px|right|alt=(Odena et al., 2016) Figure 10: Mean MS-SSIM scores for 10 ImageNet classes (y-axis) plotted against the number of classes handled by each sub-model (x-axis).|(Odena et al., 2016) Figure 10: Mean MS-SSIM scores for 10 ImageNet classes (y-axis) plotted against the number of classes handled by each sub-model (x-axis).]]<br />
<br />
=== Effect of Class Splits on Image Sample Quality ===<br />
<br />
The final experiment conducted in the paper investigates how the number of classes a fixed GAN model has to generate impacts the diversity of the generated images. In their main model the authors split the data such that each AC-GAN only had to generate images from 10 classes. Here they experimented with changing the number of classes each sub-model has to generate while holding the architecture of the sub-models fixed. They report the mean MS-SSIM score of the generated images from the original classes (when the model had only 10 to generate) at each split level. Perhaps unsurprisingly, the diversity of the outputs drops as the number of classes the model has to handle increases. Giving the model more parameters to handle the larger (and more diverse) set of classes might possibly eliminate this effect.<br />
<br />
Finally, they state that they were unable to get a regular GAN to converge on the task of generating images from one class per model. This could be due to the difficulty of training GANs and the liminted amount of training data available for each class rather than any theoretical property of class splits.<br />
<br />
= Results =<br />
<br />
=== Model ===<br />
The authors apply their AC-GAN model to the ImageNet [[#References|(Russakovsky et al., 2015)]] dataset. The data have 1000 classes which the authors split into 100 groups of 10. An AC-GAN model is trained on each group of 10 to give results reported for the paper. The authors give some examples of the images generated from this setup. They note that these are selected to show the success of the model, not give a balanced representation of how good it is:<br />
<br />
[[File:Figure_1_AC-GAN.JPG|thumb|1000px|center|alt=(Odena et al., 2016) Figure 1: Selected images generated by the AC-GAN model for the ImageNet dataset.|(Odena et al., 2016) Figure 1: Selected images generated by the AC-GAN model for the ImageNet dataset.]]<br />
<br />
The apparent value of the model is its ability to generate high-resolution samples with global coherence. From the images given above one must acknowledge they are impressively realistic. The authors link to a [https://photos.google.com/share/AF1QipPzTToH3wxrKoF8l5nWvSNz7D_oS-KB3YMQ6Ji-4XK3AtgJmlb5QCdqRQLqLSjfkw?key=YnhJb2UyMnZwb01oZ2xhUjBraE9fdWU1VVpNZTVB full set] of images generated from every ImageNet class as well. As one would expect from their acknowledgement that images above are selected to be most impressive, not all samples at the linked page are quite so coherent.<br />
<br />
[[File:Figure_9_(Bottom).JPG|thumb|400px|right|alt=(Odena et al., 2016) Figure 9 (Bottom): Each column is a different class. Each row is generated by a different latent encoding $z$.|(Odena et al., 2016) Figure 9 (Bottom): Each column is a different class. Each row is generated by a different latent encoding $z$.]]<br />
<br />
The authors compare their model with state-of-the-art results from [[#References| Salimans et al. (2016)]] on the CIFAR-10 dataset at a 32 x 32 resolution. To score the two models they use Inception Score instead of log-likelihood, which they claim is too inaccurate to be reported. Their model achieves a score of $8.25 \pm 0.07$ versus the previous state-of-the-art of $8.09 \pm 0.07$.<br />
<br />
[[#References|Odena et al. (2016)]] argue that the class conditional generator allows $G$ to learn a representation of $Z$ independent of $C$ in section 3, and give evidence of the claim later in section 4.5 by showing that images generated with a fixed latent vector $z$ but different class labels $c$ have similar global structure (e.g. orientation of the subject) but the subjects (bird species) vary according to the label. Interestingly, the background (especially in the top row) also varies with the class label. This can possibly be attributed to the bird species coming from different areas, hence a seagull might be expected to have an ocean background. Clearly here the model benefits from the fact that the authors grouped similar classes together. A more interesting analysis might show the same comparison between different classes, such as birds and forklifts, to see how global structure is encoded across them.<br />
<br />
The authors also include a discussion of whether their model is overfitting the training data. Their first test is to find the nearest neighbour in the training set of a generated image by the L1 measure in pixel space and visually compare the two images. This is a fairly naive approach, since the L1 loss in pixel space is extremely unlikely to identify whether two images are perceptually similar. Here would have been a good place to use the MS-SSIM metric to identify the nearest neighbours, since it is intended to measure perceptual similarity. The images they provide from this analysis are below.<br />
<br />
[[File:Figure_8.JPG|thumb|500px|center|alt=(Odena et al., 2016) Figure 8: Images generated by the AC-GAN model (top) and their nearest neighbour (L1 measure in pixel space) in the ImageNet dataset (bottom).|(Odena et al., 2016) Figure 8: Images generated by the AC-GAN model (top) and their nearest neighbour (L1 measure in pixel space) in the ImageNet dataset (bottom).]]<br />
<br />
A second check they make is that interpolating between two generated images in the latent space does not result in any discrete transitions or holes in the image interpolation. Such results would be indicative of overfitting. The <br />
images they give as evidence that this is not the case are below.<br />
<br />
[[File:Figure_9_(Top).JPG|thumb|1000px|center|alt=(Odena et al., 2016) Figure 9 (Top): Latent space interpolations of the AC-GAN model.|(Odena et al., 2016) Figure 9 (Top): Latent space interpolations of the AC-GAN model.]]<br />
<br />
=== Image Diversity Metric ===<br />
<br />
Another result the authors report is the performance of their image diversity metric. It is difficult to evaluate quantitatively, but visually we see that the scores do appear to capture the perceptual diversity of the generated class. For example the 'artichoke' generator appears to have collapsed, and has a high score, while the 'promontory' generator seems fairly diverse and has a low score:<br />
<br />
[[File:Figure_3.JPG|thumb|700px|center|alt=(Odena et al., 2016) Figure 3: Mean MS-SSIM scores and sample images from selected generated (top) and real (bottom) ImageNet classes.|(Odena et al., 2016) Figure 3: Mean MS-SSIM scores and sample images from selected generated (top) and real (bottom) ImageNet classes.]]<br />
<br />
= Critique =<br />
=== Model ===<br />
<br />
A major attraction of this paper is the impressive quality of samples generated by the model. GANs often generate samples that are locally plausible but globally not realistic (e.g. a generated image of a dog has fur but the overall shape is not distinguishable). As we have seen in this critique, and as acknowledged by the authors, the most impressive samples are not representative of the model's overall performance.<br />
<br />
The model itself is not a very big advancement of the field. It combines two ideas that are both already prevalent in the research without any other justification than that it seems like a natural thing to do. As [https://openreview.net/forum?id=BkDDM04Ke other reviewers] have noted, investigating how much value the proposed model adds by comparing it with other models that only implement one (or neither) of the changes would have made this paper a slightly more interesting read.<br />
<br />
Another criticism I have about the paper is about how they report their results. To compare with [[#References|Salimans et al. (2016)]] they use Inception Score rather than log-likelihood, which they claim is the standard. Even if their model performed worse by that measure it ought to be included with the caveat they mentioned. The models are evaluated on a different dataset and at a lower spatial resolution than was used for the rest of the paper. By the Inception Score their results are better on average but might not be significantly different given how close they are. Finally, they did not apply the mean MS-SSIM score they developed in this paper to evaluate their model against [[#References|Salimans et al. (2016)]]. This would have been a natural point to make, but instead they generate four samples from each model as their evidence.<br />
<br />
An analysis the authors could have included that was touched upon but not explored in section 4.6 of the paper, and in the [[#Results|Results]] section of this summary, is how the similarity of the classes grouped in each sub-model impact the quality of generated samples. The example I gave above was to compare generated images with the same latent code but very different classes, such as birds and forklifts, to see how the global structure transferred across dissimilar classes.<br />
<br />
A last point to make about the model section is that the authors make some unsupported claims in their discussions of the model's properties. Specifically, they state that their modification to the standard GAN formulation appears to stabilize training but offer no evidence. Another example is their claim that "AC-GANs learn a representation for $z$ that is independent of class label". They site [[#References|Kingma et al (2014)]] as evidence of this. From my review of that paper it does not appear that the authors gave evidence for such a claim.<br />
<br />
=== GAN Quality Metrics ===<br />
<br />
The Inception accuracy metric proposed in this paper has the drawback that it is only applicable in a conditional GAN setting, since the standard GAN framework has no ground-truth labels. It is also true that using a pre-trained classifier is only a proxy for determining how much generated images look like the class they are meant to represent, since classifiers are not perfect. Consider the phenomenon of adversarial attacks on classifiers to see this point. However the advantages the authors list, that the Inception accuracy can be computed on a per-class basis and is easier to interpret than the Inception Score do have merit. The metric does make sense for the task the authors use it for.<br />
<br />
The same can be said for the mean MS-SSIM metric developed in this paper. Visually it appears to be a good indicator of diversity in the GAN's outputs. The authors claim that the mean MS-SSIM is a fast and easy-to-compute metric for perceptual variability and collapsing behaviour in a GAN. It is unclear how fast the metric can be computed since for each class the MS-SSIM has to be computed 100*99 times, once for each pair of images. The authors do not discuss how quickly it can actually be done.<br />
<br />
=== Experimental Results on GAN Properties ===<br />
<br />
The authors included three analyses which I have termed experiments. Of these, the first concluded that images generated at a higher resolution are more dscriminable than images generated at lower resolutions, even when they have been resized to be comparable. This does not seem like a very revolutionary conclusion. For one thing, the space of lower resolution images in contained in the higher resolution space. In essence the high resolution model could generate lower resolution images by setting blocks of 4 pixels to the same intensity. It seems unsurprising then that the lower resolution is less discriminable on average. Another reason for this could be that the high resolution model has more parameters, and is trained on higher resolution data, so it has more information with which to reconstruct class information. Finally the authors give a graph of accuracies to show this property, and the average line appears compelling, however the standard errors about the lines suggest they may not be significantly different.<br />
<br />
The second experiment is on the interaction between the Inception accuracy and mean MS-SSIM metric. The author found that they are negatively correlated, and thus that classes that are high quality also tend to be diverse. This is contrary to prevailing wisdom, and since the correlation between them is weak, it appears that it may be only a fluke of the metrics.<br />
<br />
The final experiment is on the effect of class splits on image diversity. The authors found that increasing the number of classes handled by each model reduced the diversity of generated images. They make the claim at the beginning of the paper that they show the number of classes is what makes ImageNet synthesis difficult for GANs. This analysis does point in that direction but is not quite conclusive about the issue. Another analysis they could have included towards showing this is how their Inception accuracy metric and the Inception Score are affected by the number of class splits in their model.<br />
<br />
= Conclusion =<br />
<br />
This paper's main contributions were to introduce a slight variation on previous GAN models, as well as two metrics that can be used to assess the quality of generated images. The modified GAN, dubbed the Auxiliary Classifier GAN, was shown to produce high quality, high resolution samples from ImageNet, but not consistently The authors could have done more to show why their proposed architecture was an improvement over previous methods.<br />
<br />
The metrics introduced are both fairly straightforward and appear to function as they are intended. This being said, the authors could have used them more consistently throughout the paper (such as using the MS-SSIM to find nearest neighbours instead of the L1 pixel space loss). This paper was accepted to ICML 2017 but rejected by ICLR 2018 due to the incremental nature of the model development and the ad hoc nature of the other analyses presented above.<br />
<br />
= References =<br />
# Odena, A., Olah, C., & Shlens, J. (2016). Conditional image synthesis with auxiliary classifier gans. arXiv preprint [http://proceedings.mlr.press/v70/odena17a.html arXiv:1610.09585].<br />
# Goodfellow, I., Pouget-Abadie, J., Mirza, M., Xu, B., Warde-Farley, D., Ozair, S., ... & Bengio, Y. (2014). Generative adversarial nets. In Advances in neural information processing systems (pp. 2672-2680).<br />
# Salimans, T., Goodfellow, I., Zaremba, W., Cheung, V., Radford, A., & Chen, X. (2016). Improved techniques for training gans. In Advances in Neural Information Processing Systems (pp. 2234-2242).<br />
# Karras, T., Aila, T., Laine, S., & Lehtinen, J. (2017). Progressive Growing of GANs for Improved Quality, Stability, and Variation. arXiv preprint [https://arxiv.org/abs/1710.10196 arXiv:1710.10196].<br />
# Mirza, M., & Osindero, S. (2014). Conditional generative adversarial nets. arXiv preprint [https://arxiv.org/abs/1411.1784 arXiv:1411.1784].<br />
# van den Oord, A., Kalchbrenner, N., Espeholt, L., Vinyals, O., & Graves, A. (2016). Conditional image generation with pixelcnn decoders. In Advances in Neural Information Processing Systems (pp. 4790-4798).<br />
# Reed, S. E., Akata, Z., Mohan, S., Tenka, S., Schiele, B., & Lee, H. (2016). Learning what and where to draw. In Advances in Neural Information Processing Systems (pp. 217-225).<br />
# Szegedy, C., Liu, W., Jia, Y., Sermanet, P., Reed, S., Anguelov, D., ... & Rabinovich, A. (2015). Going deeper with convolutions. In Proceedings of the IEEE conference on computer vision and pattern recognition (pp. 1-9).<br />
# Sutskever, I., Vinyals, O., & Le, Q. V. (2014). Sequence to sequence learning with neural networks. In Advances in neural information processing systems (pp. 3104-3112).<br />
# Ramsundar, B., Kearnes, S., Riley, P., Webster, D., Konerding, D., & Pande, V. (2015). Massively multitask networks for drug discovery. arXiv preprint [https://arxiv.org/abs/1502.02072 arXiv:1502.02072]<br />
# Odena, A. (2016). Semi-supervised learning with generative adversarial networks. arXiv preprint [https://arxiv.org/abs/1606.01583 arXiv:1606.01583].<br />
# Theis, L., Oord, A. V. D., & Bethge, M. (2015). A note on the evaluation of generative models. arXiv preprint [https://arxiv.org/abs/1511.01844 arXiv:1511.01844].<br />
# Wang, Z., Simoncelli, E. P., & Bovik, A. C. (2003, November). Multiscale structural similarity for image quality assessment. In Signals, Systems and Computers, 2004. Conference Record of the Thirty-Seventh Asilomar Conference on (Vol. 2, pp. 1398-1402). IEEE.<br />
# Russakovsky, O., Deng, J., Su, H., Krause, J., Satheesh, S., Ma, S., ... & Berg, A. C. (2015). Imagenet large scale visual recognition challenge. International Journal of Computer Vision, 115(3), 211-252.</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Unsupervised_Domain_Adaptation_with_Residual_Transfer_Networks&diff=31114Unsupervised Domain Adaptation with Residual Transfer Networks2017-11-22T19:30:19Z<p>Nkadiyar: /* References */</p>
<hr />
<div>== Introduction ==<br />
'''Domain Adaptation''' [https://en.wikipedia.org/wiki/Domain_adaptation]is a problem in machine learning which involves taking a model which has been trained on a source domain, and applying this to a different (but related) target domain. '''Unsupervised domain adaptation''' refers to the situation in which the source data is labelled, while the target data is (predominantly) unlabeled. The problem at hand is then finding ways to generalize the learning on the source domain to the target domain. In the age of deep networks this problem has become particularly salient due to the need for vast amounts of labeled training data, in order to reap the benefits of deep learning. Manual generation of labeled data is often prohibitive, and in absence of such data networks are rarely performant. The attempt to circumvent this drought of data typically necessitates the gathering of "off-the-shelf" data sets, which are tangentially related and contain labels, and then building models in these domains. The fundamental issue that unsupervised domain adaptation attempts to address is overcoming the inherent shift in distribution across the domains, without the ability to observe this shift directly.<br />
<br />
This paper proposes a method for unsupervised domain adaptation which relies on three key components: <br />
# A kernel-based penalty to ensure that the abstract representations generated by the networks hidden layers are similar between the source and the target data; <br />
# An entropy based penalty on the target classifier, which exploits the entropy minimization principle; and <br />
# A residual network structure is appended, which allows the source and target classifiers to differ by a (learned) residual function, thus relaxing the shared classifier assumption which is traditionally made.<br />
<br />
This method outperforms state-of-the-art techniques on common benchmark datasets, and is flexible enough to be applied in most feed-forward neural networks.<br />
<br />
[[File:Source-and-Target-Domain-Office-31-Backpack.png|thumb|right|The Office-31 Dataset Images for Backpack. Shows the variation in the source and target domains to motivate why these methods are important.]] <br />
=== Working Example (Office-31) === <br />
In order to assist in the understanding of the methods, it is helpful to have a tangible sense of the problem front of mind. The Domain Adaptation Project [https://people.eecs.berkeley.edu/~jhoffman/domainadapt/] provides data sets which are tailored to the problem of unsupervised domain adaptation. One of these data sets (which is later used in the experiments of this paper) has images which are labeled based on the Amazon product page for the various items. There are then corresponding pictures taken either by webcams or digital SLR cameras. The goal of unsupervised domain adaptation on this data set would be to take any of the three image sources as the source domain, and transfer a classifier to the other domain; see the example images to understand the differences.<br />
<br />
One can imagine that, while it is likely easy to scrape labeled images from Amazon, it is likely far more difficult to collect labeled images from webcam or DSLR pictures directly. The ultimate goal of this method would be to train a model to recognize a picture of a backpack taken with a webcam, based on images of backpacks scraped from Amazon (or similar tasks).<br />
<br />
== Related Work ==<br />
Broadly speaking, the problem of domain adaptation mitigates manual labeling of data in areas such as machine learning, computer vision, and natural language processing. The general goal of domain adaptation is to reduce the discrepancy in probability distributions between the source and target domains.<br />
<br />
Research into the use of Deep Neural Networks for the purpose of domain adaptation has suggested that, while networks learn abstract feature representations which can reduce the discrepancy across domains, it is not possible to wholly remove it [http://www.icml-2011.org/papers/342_icmlpaper.pdf], [https://arxiv.org/pdf/1412.3474.pdf]. Further work has been done to design networks which adapt traditional deep nets (typically CNNs) to specifically address the problems posed by domain adaptation, these methods all only address the issue of feature adaptation [https://arxiv.org/pdf/1502.02791.pdf], [https://arxiv.org/pdf/1409.7495.pdf], [https://people.eecs.berkeley.edu/~jhoffman/papers/Tzeng_ICCV2015.pdf]. That is, they all assume that the target and source classifiers are shared between domains. <br />
<br />
The authors drew particular motivation from He et al. [https://arxiv.org/abs/1512.03385] with the proposed structure of residual networks. Combining the insights from the ResNet architecture, in addition to previous work that had leveraged classifier adaptation (in the context where some target data is labeled) [http://citeseerx.ist.psu.edu/viewdoc/download?doi=10.1.1.130.8224&rep=rep1&type=pdf], [http://www.machinelearning.org/archive/icml2009/papers/445.pdf], [http://ieeexplore.ieee.org/document/5539870/] the authors develop their proposed network.<br />
<br />
== Residual Transfer Networks ==<br />
Generally, in an unsupervised domain adaptation problem, we are dealing with a set $\mathcal{D}_s$ (called the source domain) which is defined by $\{(x_i^s, y_i^s)\}_{i=1}^{n_s}$. That is the set of all labeled input-output pairs in our source data set. We denote the number of source elements by $n_s$. There is a corresponding set $\mathcal{D}_t = \{(x_i^t)\}_{i=1}^{n_t}$ (the target domain), consisting of unlabeled input values. There are $n_t$ such values. <br />
[[File:RTN-Structure.png|thumb|left|upright|The overarching structure of the RTN. Consists of an existing network, to which a bottleneck, MMD block, and residual block is appended.]]<br />
We can think of $\mathcal{D}_s$ as being sampled from some underlying distribution $p$, and $\mathcal{D}_t$ as being sampled from $q$. Generally we have that $p \neq q$, partially motivating the need for domain adaptation methods. <br />
<br />
We can consider the classifiers $f_s(\underline{x})$ and $f_t(\underline{x})$, for the source domain and target domain respectively. It is possible to learn $f_s$ based on the sample $\mathcal{D}_s$. Under the '''shared classifier assumption''' it would be the case that $f_s(\underline{x}) = f_t(\underline{x})$, and thus learning the source classifier is enough. This method relaxes this assumption, assuming that in general $f_s \neq f_t$, and attempting to learn both.<br />
<br />
The example network extends deep convolutional networks (in this case AlexNet [http://vision.stanford.edu/teaching/cs231b_spring1415/slides/alexnet_tugce_kyunghee.pdf]) to '''Residual Transfer Networks''', the mechanics of which are outlined below. Recall that, if $L(\cdot, \cdot)$ is taken to be the cross-entropy loss function, then the empirical error of a CNN on the source domain $\mathcal{D}_s$ is given by:<br />
<br />
<center><br />
<math display="block"><br />
\min_{f_s} \frac{1}{n_s} \sum_{i=1}^{n_s} L(f_s(x_i^s), y_i^s)<br />
</math> <br />
</center><br />
<br />
In a standard implementation, the CNN optimizes over the above loss. This will be the starting point for the RTN.<br />
<br />
=== Structural Overview ===<br />
The model proposed in this paper extends existing CNN's and alters the loss function that is optimized over. While each of these components is discussed in depth below, the overarching architecture involves four components:<br />
<br />
# An existing deep model. While this can be any model, in theory, the authors leverage AlexNet in practice.<br />
# A bottleneck layer, used to reduce the dimensionality of the learned abstract feature space, directly after the existing network.<br />
# An MMD block, with the expressed intention of feature adaptation.<br />
# A residual block, with the expressed intention of classifier adaptation. <br />
<br />
This structure is then optimized against a loss function which combines the standard cross-entropy penalty with MMD and target entropy penalties, yielding the proposed Residual Transfer Network (RTN) structure.<br />
<br />
=== Feature Adaptation ===<br />
Feature adaptation refers to the process in which the features which are learned to represent the source domain are made applicable to the target domain. Broadly speaking a CNN works to generate abstract feature representations of the distribution that the inputs are sampled from. It has been found that using these deep features can reduce, but not remove, cross-domain distribution discrepancy, hence the need for feature adaptation. It is important to note that CNN's transfer from general to specific features as the network gets deeper. In this light, the discrepancy between the feature representation of the source and the target will grow through a deeper convolutional net. As such a technique for forcing these distributions to be similar is needed.<br />
<br />
In particular the authors of this paper impose a bottleneck layer (call it $fc_b$) which is included after the final convolutional layer of AlexNet. This dense layer is connected to an additional dense layer $fc_c$, (which will serve as the target classification layer). They then compute the tensor product between the activations of the layers, performing "lossless multi-layer feature fusion". That is for the source domain they define $z_i^s \overset{\underset{\mathrm{def}}{}}{=} x_i^{s,fc_b}\otimes x_i^{s,fc_c}$ and for the target domain, $z_i^t \overset{\underset{\mathrm{def}}{}}{=} x_i^{t,fc_b}\otimes x_i^{t,fc_c}$. The authors then employ feature adaptation by means of Maximum Mean Discrepancy, between the source and target domains, on these fusion features.<br />
<br />
[[File:RTN-MMD-Block.png|right|thumb|The Maximum Mean Discrepancy Block (MMD) included in the RTN. The outputs of $fc_b$ and $fc_c$ are fused through a tensor product, and then passed through the MMD penalty, ensuring distributional similarity.]]<br />
<br />
==== Maximum Mean Discrepancy ==== <br />
The Maximum Mean Discrepancy (MMD) is a Kernel method that involves mapping to a Reproducing Kernel Hilbert Space (RKHS) [https://en.wikipedia.org/wiki/Reproducing_kernel_Hilbert_space]. Denote the RKHS $\mathcal{H}_K$ with a characteristic kernel $K$. We then define the '''mean embedding''' of a distribution $p$ in $\mathcal{H}_K$ to be the unique element $\mu_K(p)$ such that $\mathbf{E}_{x\sim p}f(x) = \langle f(x), \mu_K(p)\rangle_{\mathcal{H}_K}$ for all $f \in \mathcal{H}_K$. Now, if we take $\phi: \mathcal{X} \to \mathcal{H}_K$, then we can define the MMD between two distributions $p$ and $q$ as follows:<br />
<br />
<center><br />
<math display="block"><br />
d_k(p, q) \overset{\underset{\mathrm{def}}{}}{=} ||\mathbf{E}_{x\sim p}(\phi(x^s)) - \mathbf{E}_{x\sim q}(\phi(x^t))||_{\mathcal{H}_K}<br />
</math><br />
</center><br />
<br />
Effectively, the MMD will compute the self-similarity of $p$ and $q$, and subtract twice the cross-similarity between the distributions: $\widehat{\text{MMD}}^2 = \text{mean}(K_{pp}) + \text{mean}(K_{qq}) - 2\times\text{mean}(K_{pq})$. From here we can infer that $p$ and $q$ are equivalent distributions if and only if the $\text{MMD} = 0$. If we then wish to force two distributions to be similar, this becomes a minimization problem over the MMD.<br />
<br />
Two important notes:<br />
# The RKHS, and as such MMD, depend on the choice of the kernel;<br />
# Computing the MMD efficiently requires an unbiased estimate of the MMD (as outlined [https://arxiv.org/pdf/1502.02791.pdf]).<br />
<br />
==== MMD for Feature Adaptation in the RTN ====<br />
The authors wish to minimize the MMD between the fusion features outlined above derived from the source and target domains. Concretely this amounts to forcing the distribution of the abstract representation of the source domain $\mathcal{D}_s$ to be similar to the distribution of the abstract representation of the target domain $\mathcal{D}_t$. Performing this optimization over the fused features between the $fb_b$ and $fb_c$ forces each of those layers towards similar distributions.<br />
<br />
Practically this involves an additional penalty function given by the following:<br />
<br />
<center><br />
<math display="block"><br />
D_{\mathcal{L}}(\mathcal{D}_s, \mathcal{D}_t) = \sum_{i,j=1}^{n_s} \frac{k(z_i^s, z_j^s)}{n_s^2} + \sum_{i,j=1}^{n_t} \frac{k(z_i^t, z_j^t)}{n_t^2} + \sum_{i=1}^{n_s}\sum_{j=1}^{n_t} \frac{k(z_i^s, z_j^t)}{n_sn_t} <br />
</math><br />
</center><br />
<br />
Where the characteristic kernel $k(z, z')$ is the Gaussian kernel, defined on the vectorization of tensors, with bandwidth parameter $b$. That is: $k(z, z') = \exp(-||vec(z) - vec(z')||^2/b)$.<br />
<br />
=== Classifier Adaptation ===<br />
In traditional unsupervised domain adaptation there is a '''shared-classifier assumption''' which is made. In essence, if $f_s(x)$ represents the classifier on the source domain, and $f_t(x)$ represents the classifier on the target domain then this assumption simply states that $f_s = f_t$. While this may seem to be a reasonable assumption at first glance, it is problematic largely in that this is an assumption that is incredibly difficult to check. If it could be readily confirmed that the source and target classifiers could be shared, then the problem of domain adaptation would be largely trivialized. Instead, the authors here relax this assumption slightly. They postulate that instead of being equivalent, the source and target classifier differ by some perturbation function $\Delta f$. The general idea is that, by assuming $f_S(x) = f_T(x) + \Delta f(x)$, where $f_S$ and $f_T$ correspond to the source and target classifiers, pre-activation, and $\Delta f(x)$ is some residual function.<br />
<br />
The authors then suggest using residual blocks, as popularized by the ResNet framework [https://arxiv.org/pdf/1512.03385.pdf], to learn this residual function.<br />
<br />
[[File:Residual-Block-vs-DNN.png|thumb|left|A comparison of a standard Deep Neural Network block which is designed to fit a function H(x) compared to a residual block which fits H(x) as the sum of the input, x, and a learned residual function, F(X).]]<br />
==== Residual Networks Framework ==== <br />
A (Deep) Residual Network, as proposed initially in ResNet, employs residual blocks to assist in the learning process, and were a key component of being able to train extraordinarily deep networks. The Residual Network is comprised largely in the same manner as standard neural networks, with one key difference, namely the inclusion of residual blocks - sets of layers which aim to estimate a residual function in place of estimating the function itself. <br />
<br />
That is, if we wish to use a DNN to estimate some function $h(x)$, a residual block will decompose this to $h(x) = F(x) + x$. The layers are then used to learn $F(x)$, and after the layers which aim to learn this residual function, the input $x$ is recombined through element-wise addition, to form $h(x) = F(x) + x$. This was initially proposed as a manner to allow for deeper networks to be effectively trained, but has since used in novel contexts.<br />
<br />
==== Residual Blocks in the RTN ====<br />
[[File:RTN-Residual-Block.png|thumb|right|The Structure of the Residual Block in the RTN framework. The block relies on two additional dense layers following the target classifier in an attempt to learn the residual difference between the source and target classifiers.]] The authors leverage residual blocks for the purpose of classifier adaptation. Operating under the assumption that the source and target classifiers differ by an arbitrary perturbation function, $f(x)$, the authors add an additional set of densely connected layers which the source data will flow through. In particular, the authors take the $fc_c$ layer above as the desired target classifier. For the source data an additional set of layers ($fc-1$ and $fc-2$) are added following $fc_c$, which are connected as a residual block. The output of the classifier layer is then added back to the output of the residual block in order to form the source classifier.<br />
<br />
It is necessary to note that in this case the output from $fc_c$ passes the non-activated (i.e. pre-softmax activation) to the element-wise addition, the result of which is passed through the activation layer, yielding the source prediction. In the provided diagram, we have that $f_S(x)$ represents the non-activated output from the additive layer in the residual block; $f_T(x)$ represents the non-activated output from the target classifier; and $fc-1$/$fc-2$ are used to learn the perturbation function $\Delta f(x)$.<br />
<br />
==== Entropy Minimization ====<br />
In addition to the residual blocks, the authors make use of the '''entropy minimization principle''' [http://www.iro.umontreal.ca/~lisa/pointeurs/semi-supervised-entropy-nips2004.pdf] to further refine the classifier adaptation. In particular, by minimizing the entropy of the target classifier (or more correctly, the entropy of the class conditional distribution $f_j^t(x_i^t) = p(y_i^t = j \mid x_i^t; f_t)$), low-density separation between the classes is encouraged. '''Low-Density Separation''' is a concept used predominantly in semi-supervised learning, which in essence tries to draw class decision boundaries in regions where there are few data points (labeled or unlabeled). The above paper leverages an entropy regularization scheme to achieve the goal low-density separation goal; this is adopted here to the case of unsupervised domain adaptation.<br />
<br />
In practice this amounts to adding a further penalty based on the entropy of the class conditional distribution. In particular, if $H(\cdot)$ is defined to be the entropy function, such that $H(f_t(x_i^t)) = - \sum_{j=1}^c f_j^t(x_i^t)\log f_j^t(x_i^t)$, where $c$ is the number of classes and $f_j^t(x_i^t)$ represents the probability of predicting class $j$ for point $x_i^t$, then over the target domain $\mathcal{D}_t$ we define the entropy penalty to be:<br />
<br />
<center><br />
<math display="block"><br />
\frac{1}{n_t} \sum_{i=1}^{n_t} H(f_t(x_i^t))<br />
</math><br />
</center><br />
<br />
The combination of the residual learning and the entropy penalty, the authors hypothesize will enable effective classifier adaptation.<br />
<br />
=== Residual Transfer Network ===<br />
The combination of the MMD loss introduced in feature adaptation, the residual block introduced in classifier adaptation, and the application of the entropy minimization principle cumulates in the Residual Transfer Network proposed by the authors. The model will be optimized according to the following loss function, which combines the standard cross-entropy, MMD penalty, and entropy penalty:<br />
<br />
<center><br />
<math display="block"><br />
\min_{f_s = f_t + \Delta f} \underbrace{\left(\frac{1}{n_s} \sum_{i=1}^{n_s} L(f_s(x_i^s), y_i^s)\right)}_{\text{Typical Cross-Entropy}} + \underbrace{\frac{\gamma}{n_t}\left(\sum_{i=1}^{n_t} H(f_t(x_i^t)) \right)}_{\text{Target Entropy Minimization}} + \underbrace{\lambda\left(D_{\mathcal{L}}(\mathcal{D}_s, \mathcal{D}_t)\right)}_{\text{MMD Penalty}}<br />
</math><br />
</center><br />
<br />
Where we take $\gamma$ and $\lambda$ to be tradeoff parameters between the entropy penalty and the MMD penalty. As classifier adaptation proposed in this paper and feature adaptation studied in [5, 6] are tailored to adapt different layers of deep networks, they are expected to complement each other and to establish better performance.<br />
<br />
The full network, which is trained subject to the above optimization problem, thus takes on the following structure.<br />
<br />
[[File:rtn-full-paper-structure.png||center|alt=The Structure of the RTN]]<br />
<br />
== Experiments == <br />
<br />
=== Set-up ===<br />
The performance of RTN was jointly compared across two key data sets in the area of Unsupervised Domain Adaptation. Specifically, Office-31 (discussed in the introduction) and Office-Caltech (maintained by the same project group). Office-31 is comprised of images from 3 sources, Amazon ('''A'''), Webcam ('''W'''), and DSLR ('''D'''), of 31 different objects. Office-Caltech is derived by considering 10 classes common to both the Office-31 and the Caltech data sets, thus providing further adaptation possibilities. This provides 6 Transfer Tasks on the 31 classes of Office-31 ($\{(A,W), (A,D), (W,A), (W,D), (D,A), (D,W)\}$) and 12 Transfer Tasks on the 10 classes of Office-Caltech ($\{(A,W), (A,D), (A,C), (W,A), (W,D), (W,C), (D,A), (D,W), (D,C), (C,A), (C,W), (C,D)\}$).<br />
<br />
The authors then compare the results on the 18 different adaptation tasks against 6 other models. In order to determine the efficacy of the various contributions outlined in the paper they perform an ablation study, evaluating variants of the RTN. Specifically, they consider the RTN with only the MMD module ('''RTN (mmd)'''), the RTN with the MMD module and the entropy minimization ('''RTN (mmd+ent)'''), and the complete RTN ('''RTN (mmd+ent+res)'''). The experiments leverage all the labeled training data and compute accuracy across all unlabeled domain data. The parameters of the model (i.e. $\gamma$, and $\lambda$) are fixed based on a single validation point on the transfer task $\mathbf{A}\to\mathbf{W}$. These parameters are then maintained across all transfer tasks. <br />
<br />
As for specification details, the authors use mini-batch SGD, with momentum $0.9$, and with the learning rate adjusted based on $\eta_p = \frac{\eta_0}{(1 + \alpha p)^\beta}$, where $p$ indicates the portion of training completed (linear from $0$ to $1$), $\eta_0 = 0.01$, $\alpha = 10$ and $\beta = 0.75$, which was optimized for low error on the source. The MMD and entropy parameters, set as above, were maintained at $\lambda = 0.3$ and $\gamma - 0.3$.<br />
<br />
=== Results ===<br />
[[File:table-1-results.PNG|thumb|right|Results from the Office-31 Experiment]][[File:table-2-results.PNG|thumb|right|Results from the Office-Caltech Experiment]]<br />
In aggregate, the network outperformed all comparison methods, across all transfer tasks. Broadly speaking the network saw the largest increases in accuracy on the hard transfer tasks (for instance $\mathbf{A} \to \mathbf{C}$), where the source-domain discrepancy is large. The authors take this to mean that the proposed model learns "more adaptive classifiers and transferable features for safer domain adaptation." They further indicate that standard deep learning techniques (i.e. just AlexNet) perform similarly to standard shallow techniques (TCA and GFK). Deep-transfer methods which focus on feature adaptation perform significantly better than the standard methods. The proposed RTN, which adds in additional considerations for classifier adaptation, performs even better.<br />
<br />
In addition, the ablation study found a number of interesting results:<br />
# The RTN (mmd) outperforms DAN, which is founded on a similar method, but contains multiple MMD penalties (one for each layer instead of on a bottleneck), and is as such less computationally efficient;<br />
# The addition of the entropy penalty [RTN (mmd+ent)] provides significant marginal benefit over the previous RTN (mmd);<br />
# The full RTN [RTN (mmd+ent+res)] performs the best of all variants, by diminishing returns are seen over the addition of the entropy penalty.<br />
<br />
Overall the authors claim that the RTN (mmd+ent+res) is now regarded as state-of-the-art for unsupervised domain adaptation.<br />
<br />
=== Discussion ===<br />
[[File:t-sne-embeddings.png|thumb|left|t-SNE Embeddings Comparing the Performance of DAN and RTN]] <br />
[[File:mean-sd-layer-outputs.png|thumb|right|The Mean and Standard Deviations of the outputs from the Source Classifier, Target Classifier, and Residual Functions. As expected, the residual function provides a small, but non-zero, contribution.]] <br />
[[File:gamma-tradeoff.png|thumb|left|The accuracy of tests by varying the parameter $\gamma$. We first see an increase in accuracy up to an ideal point, before having the accuracy fall again.]]<br />
[[File:classifier-shift.png|thumb|right|The corresponding weights of the classifier layers, if trained on the labeled source and target data, exhibiting the differences which exist between the two classifiers in an ideal state. ]]<br />
<br />
==== Visualizing Predictions (Versus DAN) ====<br />
DAN uses a similar method for feature adaptation but neglects any attempt at classifier adaptation (i.e. it makes the shared-classifier assumption). In order to demonstrate that this leads to the worse performance, the authors provide images showing the t-SNE embeddings by DAN and RTN on the transfer task $\mathbf{A} \to \mathbf{W}$. The images show that the target categories are not well discriminated by the source classifier, suggesting a violation of the shared-classifier assumption. Conversely, the target classifier for the RTN exhibits better discrimination.<br />
<br />
==== Layer Responses and Classifier Shift ==== <br />
The authors further consider the mean and standard deviation of the outputs of $f_S(x)$, $f_T(x)$ and $\Delta f(x)$ to consider the relative contributions of the different components. As expected, $\Delta f(x)$ provides a small (though non-zero) contribution to the learned source classifier. This provides some merit to the idea of residual learning on the classifiers. <br />
<br />
In addition, the authors train classifiers on the source and target data, with labels present, and compare the realized weights. This is used to test how different the ideal weights are on separate classifiers. The results suggest that there is, in fact, a discrepancy between the classifiers, further motivating the use of tactics to avoid the shared-classifier assumption. <br />
<br />
==== Parameter Sensitivity ==== <br />
Lastly, the authors test the sensitivity of these results against the parameter $\gamma$. They run this test on $\mathbf{A}\to\mathbf{W}$ in addition to $\mathbf{C}\to\mathbf{W}$, varying the parameter from $0.01$ to $1.0$. They find that, on both tasks, the increase of the parameter initially improves accuracy, before seeing a drop-off.<br />
<br />
== Conclusion ==<br />
This paper presented a novel approach to unsupervised domain adaptation which relaxed assumptions made by previous models with regard to the shared nature of classifiers. Like previous models this proposed network leverages feature adaptation by matching the distributions of features across the domains. In addition, using a residual network and entropy minimization tactic, the target classifier is allowed to differ from the source classifier by implementing a new residual transfer module as the bridge. In particular, this approach allows for easy integration into existing networks, and can be implemented with any standard deep learning software.<br />
<br />
For follow-up considerations, the authors propose looking for adaptations which may be useful in the semi-supervised domain adaptation problem.<br />
<br />
== Critique ==<br />
While the paper presents a clear approach, which empirically attains great results on the desired tasks, I question the benefit to the residual block that is employed. The results of the ablation study seem to suggest that the majority of the benefits can be derived from using the MMD and Entropy penalties. The residual block appears to add marginal, perhaps insignificant contributions to the outcome. Despite this, the use of MMD loss is not novel, and the entropy loss is less well documented, and less thoroughly explored. Perhaps a different set of ablations would have indicated that the three parts, indeed, are equally effective (and the diminishing returns stems from stacking the three methods), but as it is presented, I question the utility of the final structure versus a less complicated, less novel approach.<br />
<br />
==References==<br />
# https://en.wikipedia.org/wiki/Domain_adaptation<br />
# https://people.eecs.berkeley.edu/~jhoffman/domainadapt/<br />
# Glorot, Xavier, Antoine Bordes, and Yoshua Bengio. "Domain adaptation for large-scale sentiment classification: A deep learning approach." Proceedings of the 28th international conference on machine learning (ICML-11). 2011.<br />
# Tzeng, Eric, et al. "Deep domain confusion: Maximizing for domain invariance." arXiv preprint arXiv:1412.3474 (2014).<br />
# Long, Mingsheng, et al. "Learning transferable features with deep adaptation networks." International Conference on Machine Learning. 2015.<br />
# Ganin, Yaroslav, and Victor Lempitsky. "Unsupervised domain adaptation by backpropagation." International Conference on Machine Learning. 2015.<br />
# Tzeng, Eric, et al. "Simultaneous deep transfer across domains and tasks." Proceedings of the IEEE International Conference on Computer Vision. 2015.<br />
# He, Kaiming, et al. "Deep residual learning for image recognition." Proceedings of the IEEE conference on computer vision and pattern recognition. 2016.<br />
# Yang, Jun, Rong Yan, and Alexander G. Hauptmann. "Cross-domain video concept detection using adaptive svms." Proceedings of the 15th ACM international conference on Multimedia. ACM, 2007.<br />
# Duan, Lixin, et al. "Domain adaptation from multiple sources via auxiliary classifiers." Proceedings of the 26th Annual International Conference on Machine Learning. ACM, 2009.<br />
# Duan, Lixin, et al. "Visual event recognition in videos by learning from web data." IEEE Transactions on Pattern Analysis and Machine Intelligence 34.9 (2012): 1667-1680.<br />
# http://vision.stanford.edu/teaching/cs231b_spring1415/slides/alexnet_tugce_kyunghee.pdf<br />
# https://en.wikipedia.org/wiki/Reproducing_kernel_Hilbert_space<br />
#Long, Mingsheng, et al. "Learning transferable features with deep adaptation networks." International Conference on Machine Learning. 2015.<br />
#He, Kaiming, et al. "Deep residual learning for image recognition." Proceedings of the IEEE conference on computer vision and pattern recognition. 2016.<br />
# Grandvalet, Yves, and Yoshua Bengio. "Semi-supervised learning by entropy minimization." Advances in neural information processing systems. 2005.<br />
# More information on residual functions https://www.youtube.com/watch?v=urAp0DibYlY <br />
<br />
Expert review from the NIPS community can be found in https://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/99.html.<br />
<br />
Implementation Example: https://github.com/thuml/Xlearn</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Model-Agnostic_Meta-Learning_for_Fast_Adaptation_of_Deep_Networks&diff=31113Model-Agnostic Meta-Learning for Fast Adaptation of Deep Networks2017-11-22T19:05:22Z<p>Nkadiyar: /* Introduction & Background */</p>
<hr />
<div>='''Introduction & Background'''=<br />
Learning quickly is a hallmark of human intelligence, whether it involves recognizing objects from a few examples or quickly learning new skills after just minutes of experience. In this work, we propose a meta-learning algorithm that is general and model-agnostic, in the sense that it can be directly applied to any learning problem and model that is trained with a gradient descent procedure. Our focus is on deep neural network models, but we illustrate how our approach can easily handle different architectures and different problem settings, including classification, regression, and policy gradient reinforcement learning, with minimal modification. Unlike prior meta-learning methods that learn an update function or learning rule (Schmidhuber, 1987; Bengio et al., 1992; Andrychowicz et al., 2016; Ravi & Larochelle, 2017), this algorithm does not expand the number of learned parameters nor place constraints on the model architecture (e.g. by requiring a recurrent model (Santoro et al., 2016) or a Siamese network (Koch, 2015)), and it can be readily combined with fully connected, convolutional, or recurrent neural networks. It can also be used with a variety of loss functions, including differentiable supervised losses and nondifferentiable reinforcement learning objectives.<br />
<br />
The primary contribution of this work is a simple model and task-agnostic algorithm for meta-learning that trains a model’s parameters such that a small number of gradient updates will lead to fast learning on a new task. The paper shows the effectiveness of the proposed algorithm in different domains, including classification, regression, and reinforcement learning problems.<br />
<br />
==Key Idea==<br />
The key idea underlying this method is to train the model’s initial parameters such that the model has maximal performance on a new task after the parameters have been updated through one or more gradient steps computed with a small amount of data from that new task. This can be viewed from a feature learning standpoint as building an internal representation that is broadly suitable for many tasks. If the internal representation is suitable to many tasks, simply fine-tuning the parameters slightly (e.g. by primarily modifying the top layer weights in a feedforward model) can produce good results.<br />
<br />
='''Model-Agnostic Meta Learning (MAML)'''=<br />
The goal of the proposed model is rapid adaptation. This setting is usually formalized as few-shot learning.<br />
<br />
=== Problem set-up ===<br />
The goal of few-shot meta-learning is to train a model that can quickly adapt to a new task using only a few data points and training iterations. To do so, the model is trained during a meta-learning phase on a set of tasks, such that it can then be adapted to a new task using only a small number of parameter updates. In effect, the meta-learning problem treats entire tasks as training examples. <br />
<br />
Let us consider a model denoted by $f$, that maps the observation $\mathbf{x}$ into the output variable $a$. During meta-learning, the model is trained to be able to adapt to a large or infinite number of tasks. <br />
<br />
Let us consider a generic notion of task as below. Each task $\mathcal{T} = \{\mathcal{L}(\mathbf{x}_1.a_1,\mathbf{x}_2,a_2,..., \mathbf{x}_H,a_H), q(\mathbf{x}_1),q(\mathbf{x}_{t+1}|\mathbf{x}_t,a_t),H \}$, consists of a loss function $\mathcal{L}$, a distribution over initial observations $q(\mathbf{x}_1)$, a transition distribution $q(\mathbf{x}_{t+1}|\mathbf{x}_t)$, and an episode length $H$. In i.i.d. supervised learning problems,<br />
the length $H =1$. The model may generate samples of length $H$ by choosing an output at at each time $t$. The cost $\mathcal{L}$ provides a task-specific feedback, which is defined based on the nature of the problem. <br />
<br />
A distribution over tasks is denoted by $p(\mathcal{T})$. In the K-shot learning setting, the model is trained to learn a new task $\mathcal{T}_i$ drawn from $p(\mathcal{T})$ from only K samples drawn from $q_i$ and feedback $\mathcal{L}_{\mathcal{T}_i}$ generated by $\mathcal{T}_i$. During meta-training, a task $\mathcal{T}_i$ is sampled from $p(\mathcal{T})$, the model is trained with K samples and feedback from the corresponding loss LTi from Ti, and then tested on new samples from Ti. The model f is then improved by considering how the test error on new data from $q_i$ changes with respect to the parameters. In effect, the test error on sampled tasks $\mathcal{T}_i$ serves as the training error of the meta-learning process. At the end of meta-training, new tasks are sampled from $p(\mathcal{T})$, and meta-performance is measured by the model’s performance after learning from K samples.<br />
<br />
=== MAML Algorithm ===<br />
[[File:model.png|200px|right|thumb|Figure 1: Diagram of the MAML algorithm]]<br />
The paper proposes a method that can learn the parameters of any standard model via meta-learning in such a way as to prepare that model for fast adaptation. The intuition behind this approach is that some internal representations are more transferrable than others. Since the model will be fine-tuned using a gradient-based learning rule on a new task, we will aim to learn a model in such a way that this gradient-based learning rule can make rapid progress on new tasks drawn from $p(\mathcal{T})$, without overfitting. In effect, we will aim to find model parameters that are sensitive to changes in the task, such that small changes in the parameters will produce large improvements on the loss function of any task drawn from $p(\mathcal{T})$, see Fig 1.<br />
<br />
Note that there is no assumption about the form of the model. Only assumption is that it is parameterized by a vector of parameters $\theta$, and the loss is smooth so that the parameters can be leaned using gradient-based techniques. Formally lets assume that the model is denoted by $f_{\theta}$. When adapting<br />
to a new task $\mathcal{T}_i $, the model’s parameters $\theta$ become $\theta_i'$. In our method, the updated parameter vector $\theta_i'$ is computed using one or more gradient descent updates on task $\mathcal{T}_i $. For example, when using one gradient update:<br />
<br />
$$<br />
\theta_i ' = \theta - \alpha \nabla_{\theta \mathcal{L}_{\mathcal{T}_i}}(f_{\theta}).<br />
$$<br />
<br />
Here $\alpha$ is a the learning rate of each task and considered as a hyperparameter. They consider a single step of update for the rest of the paper, for the sake of the simplicity. <br />
<br />
The model parameters are trained by optimizing for the performance<br />
of $f_{\theta_i'}$ with respect to $\theta$ across tasks sampled from $p(\mathcal{T})$. More concretely, the meta-objective is as follows: <br />
<br />
$$<br />
\min_{\theta} \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta_i'}) = \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta - \alpha \nabla_{\theta \mathcal{L}_{\mathcal{T}_i}}(f_{\theta})})<br />
$$<br />
<br />
Note that the meta-optimization is performed over the model parameters $\theta$, whereas the objective is computed using the updated model parameters $\theta'$. The model aims to optimize the model parameters such that one or a small number of gradient step on a new task will produce maximally effective behavior on that task. <br />
<br />
Therefore the meta-learning across the tasks is performed via stochastic gradient descent (SGD), such that the model parameters $\theta$ are updated as:<br />
<br />
$$<br />
\theta \gets \theta - \beta \nabla_{\theta } \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta_i'})<br />
$$<br />
where $\beta$ is the meta step size. Outline of the algorithm is shown in Algorithm 1. <br />
[[File:ershad_alg1.png|500px|center|thumb]]<br />
<br />
The MAML meta-gradient update involves a gradient through a gradient. Computationally, this requires an additional backward pass through f to compute Hessian-vector products, which is supported by standard deep learning libraries such as TensorFlow.<br />
<br />
='''Different Types of MAML'''=<br />
In this section the MAML algorithm is discussed for different supervised learning and reinforcement learning tasks. The differences between each of these tasks are in their loss function and the way the data is generated. In general, this method does not require additional model parameters nor using any additional meta-learner to learn the update of parameters. Compared to other approaches that tend to “learn to compare new examples in a learned metric space using e.g. Siamese networks or recurrence with attention mechanisms”, the proposed method can be generalized to any other problems including classification, regression and reinforcement learning. <br />
<br />
=== Supervised Regression and Classification ===<br />
Few-shot learning is well-studied in this field. For these two types of tasks the horizon $H$ is equal to 1, since the data points are generated i.i.d. <br />
<br />
Although any common classification and regression objectives can be used as the loss, the paper uses the following losses for these two tasks. <br />
<br />
Regression : For regression we use the mean-square error (MSE):<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i} (f_{\theta}) = \sum \limits_{\mathbf{x}^{(j)}, \mathcal{y}^{(j)} \sim \mathcal{T}_i} \parallel f_{\theta} (\mathbf{x}^{(j)} - \mathbf{y}^{(j)})\parallel_2^2<br />
$$<br />
<br />
where $\mathbf{x}^{(j)}$ and $\mathbf{y}^{()j}$ are the input/output pair sampled from task $\mathcal{T}_i$. In K-shot regression tasks, K input/output pairs are provided for learning for each task. <br />
<br />
Classification: For classification we use the cross entropy loss:<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i} (f_{\theta}) = \sum \limits_{\mathbf{x}^{(j)}, \mathcal{y}^{(j)} \sim \mathcal{T}_i} \mathbf{y}^{(j)} \log f_{\theta}(\mathbf{x}^{(j)}) + (1-\mathbf{y}^{(j)}) \log (1-f_{\theta}(\mathbf{x}^{(j)}))<br />
$$<br />
<br />
According to the conventional terminology, K-shot classification tasks use K input/output pairs from each class, for a total of $NK$ data points for N-way classification.<br />
<br />
Given a distribution over tasks, these loss functions can be directly inserted into the equations in the previous section to perform meta-learning, as detailed in Algorithm 2.<br />
[[File:ershad_alg2.png|500px|center|thumb]]<br />
<br />
=== Reinforcement Learning ===<br />
In reinforcement learning (RL), the goal of few-shot meta learning is to enable an agent to quickly acquire a policy for a new test task using only a small amount of experience in the test setting. A new task might involve achieving a new goal or succeeding on a previously trained goal in a new environment. For example an agent may learn how to navigate mazes very quickly so that, when faced with a new maze, it can determine how to reliably reach the exit with only a few samples.<br />
<br />
Each RL task $\mathcal{T}_i$ contains an initial state distribution $q_i(\mathbf{x}_1)$ and a transition distribution $q_i(\mathbf{x}_{t+1}|\mathbf{x}_t,a_t)$ , and the loss $\mathcal{L}_{\mathcal{T}_i}$ corresponds to the (negative) reward function $R$. The entire task is therefore a Markov decision process (MDP) with horizon H, where the learner is allowed to query a limited number of sample trajectories for few-shot learning. Any aspect of the MDP may change across tasks in $p(\mathcal{T})$. The model being learned, $f_{\theta}$, is a policy that maps from states $\mathbf{x}_t$ to a distribution over actions $a_t$ at each timestep $t \in \{1,...,H\}$. The loss for task $\mathcal{T}_i$ and model $f_{\theta}$ takes the form<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i}(f_{\theta}) = -\mathbb{E}_{\mathbf{x}_t,a_t \sim f_{\theta},q_{\mathcal{T}_i}} \big [\sum_{t=1}^H R_i(\mathbf{x}_t,a_t)\big ]<br />
$$<br />
<br />
<br />
In K-shot reinforcement learning, K rollouts from $f_{\theta}$ and task $\mathcal{T}_i$, $(\mathbf{x}_1,a_1,...,\mathbf{x}_H)$, and the corresponding rewards $ R(\mathbf{x}_t,a_t)$, may be used for adaptation on a new task $\mathcal{T}_i$.<br />
<br />
Since the expected reward is generally not differentiable due to unknown dynamics, we use policy gradient methods to estimate the gradient both for the model gradient update(s) and the meta-optimization. Since policy gradients are an on-policy algorithm, each additional gradient step during the adaptation of $f_{\theta}$ requires new samples from the current policy $f_{\theta_i}$ . We detail the algorithm in Algorithm 3, which has the same structure as Algorithm 2 but also which requires sampling trajectories from the environment corresponding to task $\mathcal{T}_i$ in steps 5 and 8.<br />
[[File:ershad_alg3.png|500px|center|thumb]]<br />
<br />
='''Experiments'''=<br />
<br />
=== Regression ===<br />
We start with a simple regression problem that illustrates the basic principles of MAML. Each task involves regressing from the input to the output of a sine wave, where the amplitude and phase of the sinusoid are varied between tasks. Thus, $p(\mathcal{T})$ is continuous, and the input and output both have a dimensionality of 1. During training and testing, datapoints are sampled uniformly. The loss is the mean-squared error between the prediction and true value. The regressor is a neural network model with 2 hidden layers of size 40 with ReLU nonlinearities. When training with MAML, we use one gradient update with K = 10 examples with a fixed step size 0.01, and use Adam as the metaoptimizer [2]. The baselines are likewise trained with Adam. To evaluate performance, we fine-tune a single meta-learned model on varying numbers of K examples, and compare performance to two baselines: (a) pre-training on all of the tasks, which entails training a network to regress to random sinusoid functions and then, at test-time, fine-tuning with gradient descent on the K provided points, using an automatically tuned step size, and (b) an oracle which receives the true amplitude and phase as input.<br />
<br />
We evaluate performance by fine-tuning the model learned by MAML and the pre-trained model on $K = \{ 5,10,20 \}$ datapoints. During fine-tuning, each gradient step is computed using the same $K$ datapoints. Results are shown in Fig 2.<br />
<br />
<br />
[[File:ershad_results1.png|500px|center|thumb|Figure 2: Few-shot adaptation for the simple regression task. Left: Note that MAML is able to estimate parts of the curve where there are no datapoints, indicating that the model has learned about the periodic structure of sine waves. Right: Fine-tuning of a model pre-trained on the same distribution of tasks without MAML, with a tuned step size. Due to the often contradictory outputs on the pre-training tasks, this model is unable to recover a suitable representation and fails to extrapolate from the small number of test-time samples.]]<br />
<br />
=== Classification ===<br />
<br />
For classification evaluation, Omniglot and MiniImagenet datasets are used. The Omniglot dataset consists of 20 instances of 1623 characters from 50 different alphabets. <br />
<br />
The experiment involves fast learning of N-way classification with 1 or 5 shots. The problem of N-way classification is set up as follows: select N unseen classes, provide the model with K different instances of each of the N classes, and evaluate the model’s ability to classify new instances within the N classes. For Omniglot, 1200 characters are selected randomly for training, irrespective of alphabet, and use the remaining for testing. The Omniglot dataset is augmented with rotations by multiples of 90 degrees.<br />
<br />
The model follows the same architecture as the embedding function that has 4 modules with a 3-by-3 convolutions and 64 filters, followed by batch normalization, a ReLU nonlinearity, and 2-by-2 max-pooling. The Omniglot images are downsampled to 28-by-28, so the dimensionality of the last hidden layer is 64. The last layer is fed into a softmax. For Omniglot, strided convolutions is used instead of max-pooling. For MiniImagenet, 32 filters per layer are used to reduce overfitting. In order to also provide a fair comparison against memory-augmented neural networks [3] and to test the flexibility of MAML, the results for a non-convolutional network are also provided. <br />
<br />
[[File:ershad_results2.png|500px|center|thumb|Table 1: Few-shot classification on held-out Omniglot characters (top) and the MiniImagenet test set (bottom). MAML achieves results that are comparable to or outperform state-of-the-art convolutional and recurrent models. Siamese nets, matching nets, and the memory module approaches are all specific to classification, and are not directly applicable to regression or RL scenarios. The $\pm$ shows 95% confidence intervals over tasks. ]]<br />
<br />
=== Reinforcement Learning ===<br />
Several simulated continuous control environments are used for RL evaluation. In all of the domain, the MAML model is a neural network policy with two hidden layers of size 100, and ReLU activations. The gradient updates are computed using vanilla policy gradient and trust-region policy optimization (TRPO) is used as the meta-optimizer.<br />
<br />
In order to avoid computing third derivatives, finite differences are computed to <br />
compute the Hessian-vector products for TRPO. For both learning and meta-learning updates, we use the standard linear feature baseline proposed by [4], which is fitted separately at each iteration for each sampled task in the batch. <br />
<br />
Three baseline models for the comparison are: <br />
(a) pretraining one policy on all of the tasks and then fine-tuning<br />
(b) training a policy from randomly initialized weights<br />
(c) an oracle policy which receives the parameters of the task as input, which for the tasks below corresponds to a goal position, goal direction, or goal velocity for the agent. <br />
<br />
The baseline models of (a) and (b) are fine-tuned with gradient descent with a manually tuned step size.<br />
<br />
2D Navigation: In the first meta-RL experiment, the authors study a set of tasks where a point agent must move to different goal positions in 2D, randomly chosen for each task within a unit square. The observation is the current 2D position, and actions correspond to velocity commands clipped to be in the range [-0.1; 0.1]. The reward is the negative squared distance to the goal, and episodes terminate when the agent is within 0:01 of the goal or at the horizon ofH = 100. The policy was trained with MAML <br />
to maximize performance after 1 policy gradient update using 20 trajectories. They compare adaptation to a new task with up to 4 gradient updates, each with 40 samples. Results are shown in Fig. 3.<br />
<br />
[[File:ershad_results3.png|500px|center|thumb|Figure 3: Top: quantitative results from 2D navigation task, Bottom: qualitative comparison between model learned with MAML and with fine-tuning from a pretrained network ]]<br />
<br />
Locomotion. To study how well MAML can scale to more complex deep RL problems, we also study adaptation on high-dimensional locomotion tasks with the MuJoCo simulator [5]. The tasks require two simulated robots – a planar cheetah and a 3D quadruped (the “ant”) – to run in a particular direction or at a particular velocity. In the goal velocity experiments, the reward is the negative absolute value between the current velocity of the agent and a goal, which is chosen uniformly at random between 0 and 2 for the cheetah and between 0 and 3 for the ant. In the goal direction experiments, the reward is the magnitude of the velocity in either the forward or backward direction, chosen at random for each task in p(T ). The horizon is H = 200, with 20 rollouts per gradient step for all problems except the ant forward/backward task, which used 40 rollouts per step. The results in Figure 5 show that MAML learns a model that can quickly adapt its velocity and direction with even <br />
just a single gradient update, and continues to improve with more gradient steps. The results also show that, on these challenging tasks, the MAML initialization substantially outperforms random initialization and pretraining.<br />
[[File:ershad_results4.png|500px|center|thumb|Figure 4: Reinforcement learning results for the half-cheetah and ant locomotion tasks, with the tasks shown on the far right. ]]<br />
<br />
A conceptual method to achieve fast adaptation in language modeling tasks ( not been experimented on by the authors) would be to explore methods of attaching an Attention Kernel which results in a simple and differentiable loss. It has been implemented in One-Shot Language Modeling along with state-of-the-art improvements in one-shot learning on Imagenet and Omniglot [7].<br />
<br />
='''Conclusion'''=<br />
<br />
The paper introduced a meta-learning method based on learning easily adaptable model parameters through gradient descent. The approach has a number of benefits. It is simple and does not introduce any learned parameters for meta-learning. It can be combined with any model representation that is amenable to gradient-based training, and any differentiable objective, including classification, regression, and reinforcement learning. Lastly, since the method merely produces a weight initialization, adaptation can be performed with any amount of data and any number of gradient steps, though it demonstrates state-of-the-art results on classification with only one or five examples per class. The authors also show that the method can adapt an RL agent using policy gradients and a very modest amount of experience.<br />
<br />
='''Critique'''=<br />
From my opinion, the Model-Agnostic Meta-Learning looks like a simplified curriculum learning. It treats all tasks the same over the whole training history, and does not consider the difficulty of the tasks and the adaption of neural network to the task. Curriculum learning would be a good idea to speed up the training.<br />
<br />
<br />
='''References'''=<br />
# Schmidhuber, J¨urgen. Learning to control fast-weight memories: An alternative to dynamic recurrent networks. Neural Computation, 1992.<br />
# Lake, Brenden M, Salakhutdinov, Ruslan, Gross, Jason, and Tenenbaum, Joshua B. One shot learning of simple visual concepts. In Conference of the Cognitive Science Society (CogSci), 2011.<br />
# Santoro, Adam, Bartunov, Sergey, Botvinick, Matthew, Wierstra, Daan, and Lillicrap, Timothy. Meta-learning with memory-augmented neural networks. In International Conference on Machine Learning (ICML), 2016.<br />
# Duan, Yan, Chen, Xi, Houthooft, Rein, Schulman, John, and Abbeel, Pieter. Benchmarking deep reinforcement learning for continuous control. In International Conference on Machine Learning (ICML), 2016.<br />
# Todorov, Emanuel, Erez, Tom, and Tassa, Yuval. Mujoco: A physics engine for model-based control. In International Conference on Intelligent Robots and Systems (IROS), 2012.<br />
# Videos the learned policies can be found in https://sites.google.com/view/maml.<br />
# Oriol Vinyals, Charles Blundell, Timothy Lillicrap, Koray Kavukcuoglu, Daan Wierstra. "Matching Networks for One Shot Learning". arXiv:1606.04080 [cs.LG]<br />
<br />
Implementation Example: https://github.com/cbfinn/maml</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Model-Agnostic_Meta-Learning_for_Fast_Adaptation_of_Deep_Networks&diff=31112Model-Agnostic Meta-Learning for Fast Adaptation of Deep Networks2017-11-22T18:51:41Z<p>Nkadiyar: /* Problem set-up */</p>
<hr />
<div>='''Introduction & Background'''=<br />
Learning quickly is a hallmark of human intelligence, whether it involves recognizing objects from a few examples or quickly learning new skills after just minutes of experience. In this work, we propose a meta-learning algorithm that is general and model-agnostic, in the sense that it can be directly applied to any learning problem and model that is trained with a gradient descent procedure. Our focus is on deep neural network models, but we illustrate how our approach can easily handle different architectures and different problem settings, including classification, regression, and policy gradient reinforcement learning, with minimal modification. Unlike prior meta-learning methods that learn an update function or learning rule (Schmidhuber, 1987; Bengio et al., 1992; Andrychowicz et al., 2016; Ravi & Larochelle, 2017), this algorithm does not expand the number of learned parameters nor place constraints on the model architecture (e.g. by requiring a recurrent model (Santoro et al., 2016) or a Siamese network (Koch, 2015)), and it can be readily combined with fully connected, convolutional, or recurrent neural networks. It can also be used with a variety of loss functions, including differentiable supervised losses and nondifferentiable reinforcement learning objectives.<br />
<br />
The primary contribution of this work is a simple model and task-agnostic algorithm for meta-learning that trains a model’s parameters such that a small number of gradient updates will lead to fast learning on a new task. The paper shows the effectiveness of the proposed algorithm in different domains, including classification, regression, and reinforcement learning problems.<br />
<br />
='''Model-Agnostic Meta Learning (MAML)'''=<br />
The goal of the proposed model is rapid adaptation. This setting is usually formalized as few-shot learning.<br />
<br />
=== Problem set-up ===<br />
The goal of few-shot meta-learning is to train a model that can quickly adapt to a new task using only a few data points and training iterations. To do so, the model is trained during a meta-learning phase on a set of tasks, such that it can then be adapted to a new task using only a small number of parameter updates. In effect, the meta-learning problem treats entire tasks as training examples. <br />
<br />
Let us consider a model denoted by $f$, that maps the observation $\mathbf{x}$ into the output variable $a$. During meta-learning, the model is trained to be able to adapt to a large or infinite number of tasks. <br />
<br />
Let us consider a generic notion of task as below. Each task $\mathcal{T} = \{\mathcal{L}(\mathbf{x}_1.a_1,\mathbf{x}_2,a_2,..., \mathbf{x}_H,a_H), q(\mathbf{x}_1),q(\mathbf{x}_{t+1}|\mathbf{x}_t,a_t),H \}$, consists of a loss function $\mathcal{L}$, a distribution over initial observations $q(\mathbf{x}_1)$, a transition distribution $q(\mathbf{x}_{t+1}|\mathbf{x}_t)$, and an episode length $H$. In i.i.d. supervised learning problems,<br />
the length $H =1$. The model may generate samples of length $H$ by choosing an output at at each time $t$. The cost $\mathcal{L}$ provides a task-specific feedback, which is defined based on the nature of the problem. <br />
<br />
A distribution over tasks is denoted by $p(\mathcal{T})$. In the K-shot learning setting, the model is trained to learn a new task $\mathcal{T}_i$ drawn from $p(\mathcal{T})$ from only K samples drawn from $q_i$ and feedback $\mathcal{L}_{\mathcal{T}_i}$ generated by $\mathcal{T}_i$. During meta-training, a task $\mathcal{T}_i$ is sampled from $p(\mathcal{T})$, the model is trained with K samples and feedback from the corresponding loss LTi from Ti, and then tested on new samples from Ti. The model f is then improved by considering how the test error on new data from $q_i$ changes with respect to the parameters. In effect, the test error on sampled tasks $\mathcal{T}_i$ serves as the training error of the meta-learning process. At the end of meta-training, new tasks are sampled from $p(\mathcal{T})$, and meta-performance is measured by the model’s performance after learning from K samples.<br />
<br />
=== MAML Algorithm ===<br />
[[File:model.png|200px|right|thumb|Figure 1: Diagram of the MAML algorithm]]<br />
The paper proposes a method that can learn the parameters of any standard model via meta-learning in such a way as to prepare that model for fast adaptation. The intuition behind this approach is that some internal representations are more transferrable than others. Since the model will be fine-tuned using a gradient-based learning rule on a new task, we will aim to learn a model in such a way that this gradient-based learning rule can make rapid progress on new tasks drawn from $p(\mathcal{T})$, without overfitting. In effect, we will aim to find model parameters that are sensitive to changes in the task, such that small changes in the parameters will produce large improvements on the loss function of any task drawn from $p(\mathcal{T})$, see Fig 1.<br />
<br />
Note that there is no assumption about the form of the model. Only assumption is that it is parameterized by a vector of parameters $\theta$, and the loss is smooth so that the parameters can be leaned using gradient-based techniques. Formally lets assume that the model is denoted by $f_{\theta}$. When adapting<br />
to a new task $\mathcal{T}_i $, the model’s parameters $\theta$ become $\theta_i'$. In our method, the updated parameter vector $\theta_i'$ is computed using one or more gradient descent updates on task $\mathcal{T}_i $. For example, when using one gradient update:<br />
<br />
$$<br />
\theta_i ' = \theta - \alpha \nabla_{\theta \mathcal{L}_{\mathcal{T}_i}}(f_{\theta}).<br />
$$<br />
<br />
Here $\alpha$ is a the learning rate of each task and considered as a hyperparameter. They consider a single step of update for the rest of the paper, for the sake of the simplicity. <br />
<br />
The model parameters are trained by optimizing for the performance<br />
of $f_{\theta_i'}$ with respect to $\theta$ across tasks sampled from $p(\mathcal{T})$. More concretely, the meta-objective is as follows: <br />
<br />
$$<br />
\min_{\theta} \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta_i'}) = \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta - \alpha \nabla_{\theta \mathcal{L}_{\mathcal{T}_i}}(f_{\theta})})<br />
$$<br />
<br />
Note that the meta-optimization is performed over the model parameters $\theta$, whereas the objective is computed using the updated model parameters $\theta'$. The model aims to optimize the model parameters such that one or a small number of gradient step on a new task will produce maximally effective behavior on that task. <br />
<br />
Therefore the meta-learning across the tasks is performed via stochastic gradient descent (SGD), such that the model parameters $\theta$ are updated as:<br />
<br />
$$<br />
\theta \gets \theta - \beta \nabla_{\theta } \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta_i'})<br />
$$<br />
where $\beta$ is the meta step size. Outline of the algorithm is shown in Algorithm 1. <br />
[[File:ershad_alg1.png|500px|center|thumb]]<br />
<br />
The MAML meta-gradient update involves a gradient through a gradient. Computationally, this requires an additional backward pass through f to compute Hessian-vector products, which is supported by standard deep learning libraries such as TensorFlow.<br />
<br />
='''Different Types of MAML'''=<br />
In this section the MAML algorithm is discussed for different supervised learning and reinforcement learning tasks. The differences between each of these tasks are in their loss function and the way the data is generated. In general, this method does not require additional model parameters nor using any additional meta-learner to learn the update of parameters. Compared to other approaches that tend to “learn to compare new examples in a learned metric space using e.g. Siamese networks or recurrence with attention mechanisms”, the proposed method can be generalized to any other problems including classification, regression and reinforcement learning. <br />
<br />
=== Supervised Regression and Classification ===<br />
Few-shot learning is well-studied in this field. For these two types of tasks the horizon $H$ is equal to 1, since the data points are generated i.i.d. <br />
<br />
Although any common classification and regression objectives can be used as the loss, the paper uses the following losses for these two tasks. <br />
<br />
Regression : For regression we use the mean-square error (MSE):<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i} (f_{\theta}) = \sum \limits_{\mathbf{x}^{(j)}, \mathcal{y}^{(j)} \sim \mathcal{T}_i} \parallel f_{\theta} (\mathbf{x}^{(j)} - \mathbf{y}^{(j)})\parallel_2^2<br />
$$<br />
<br />
where $\mathbf{x}^{(j)}$ and $\mathbf{y}^{()j}$ are the input/output pair sampled from task $\mathcal{T}_i$. In K-shot regression tasks, K input/output pairs are provided for learning for each task. <br />
<br />
Classification: For classification we use the cross entropy loss:<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i} (f_{\theta}) = \sum \limits_{\mathbf{x}^{(j)}, \mathcal{y}^{(j)} \sim \mathcal{T}_i} \mathbf{y}^{(j)} \log f_{\theta}(\mathbf{x}^{(j)}) + (1-\mathbf{y}^{(j)}) \log (1-f_{\theta}(\mathbf{x}^{(j)}))<br />
$$<br />
<br />
According to the conventional terminology, K-shot classification tasks use K input/output pairs from each class, for a total of $NK$ data points for N-way classification.<br />
<br />
Given a distribution over tasks, these loss functions can be directly inserted into the equations in the previous section to perform meta-learning, as detailed in Algorithm 2.<br />
[[File:ershad_alg2.png|500px|center|thumb]]<br />
<br />
=== Reinforcement Learning ===<br />
In reinforcement learning (RL), the goal of few-shot meta learning is to enable an agent to quickly acquire a policy for a new test task using only a small amount of experience in the test setting. A new task might involve achieving a new goal or succeeding on a previously trained goal in a new environment. For example an agent may learn how to navigate mazes very quickly so that, when faced with a new maze, it can determine how to reliably reach the exit with only a few samples.<br />
<br />
Each RL task $\mathcal{T}_i$ contains an initial state distribution $q_i(\mathbf{x}_1)$ and a transition distribution $q_i(\mathbf{x}_{t+1}|\mathbf{x}_t,a_t)$ , and the loss $\mathcal{L}_{\mathcal{T}_i}$ corresponds to the (negative) reward function $R$. The entire task is therefore a Markov decision process (MDP) with horizon H, where the learner is allowed to query a limited number of sample trajectories for few-shot learning. Any aspect of the MDP may change across tasks in $p(\mathcal{T})$. The model being learned, $f_{\theta}$, is a policy that maps from states $\mathbf{x}_t$ to a distribution over actions $a_t$ at each timestep $t \in \{1,...,H\}$. The loss for task $\mathcal{T}_i$ and model $f_{\theta}$ takes the form<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i}(f_{\theta}) = -\mathbb{E}_{\mathbf{x}_t,a_t \sim f_{\theta},q_{\mathcal{T}_i}} \big [\sum_{t=1}^H R_i(\mathbf{x}_t,a_t)\big ]<br />
$$<br />
<br />
<br />
In K-shot reinforcement learning, K rollouts from $f_{\theta}$ and task $\mathcal{T}_i$, $(\mathbf{x}_1,a_1,...,\mathbf{x}_H)$, and the corresponding rewards $ R(\mathbf{x}_t,a_t)$, may be used for adaptation on a new task $\mathcal{T}_i$.<br />
<br />
Since the expected reward is generally not differentiable due to unknown dynamics, we use policy gradient methods to estimate the gradient both for the model gradient update(s) and the meta-optimization. Since policy gradients are an on-policy algorithm, each additional gradient step during the adaptation of $f_{\theta}$ requires new samples from the current policy $f_{\theta_i}$ . We detail the algorithm in Algorithm 3, which has the same structure as Algorithm 2 but also which requires sampling trajectories from the environment corresponding to task $\mathcal{T}_i$ in steps 5 and 8.<br />
[[File:ershad_alg3.png|500px|center|thumb]]<br />
<br />
='''Experiments'''=<br />
<br />
=== Regression ===<br />
We start with a simple regression problem that illustrates the basic principles of MAML. Each task involves regressing from the input to the output of a sine wave, where the amplitude and phase of the sinusoid are varied between tasks. Thus, $p(\mathcal{T})$ is continuous, and the input and output both have a dimensionality of 1. During training and testing, datapoints are sampled uniformly. The loss is the mean-squared error between the prediction and true value. The regressor is a neural network model with 2 hidden layers of size 40 with ReLU nonlinearities. When training with MAML, we use one gradient update with K = 10 examples with a fixed step size 0.01, and use Adam as the metaoptimizer [2]. The baselines are likewise trained with Adam. To evaluate performance, we fine-tune a single meta-learned model on varying numbers of K examples, and compare performance to two baselines: (a) pre-training on all of the tasks, which entails training a network to regress to random sinusoid functions and then, at test-time, fine-tuning with gradient descent on the K provided points, using an automatically tuned step size, and (b) an oracle which receives the true amplitude and phase as input.<br />
<br />
We evaluate performance by fine-tuning the model learned by MAML and the pre-trained model on $K = \{ 5,10,20 \}$ datapoints. During fine-tuning, each gradient step is computed using the same $K$ datapoints. Results are shown in Fig 2.<br />
<br />
<br />
[[File:ershad_results1.png|500px|center|thumb|Figure 2: Few-shot adaptation for the simple regression task. Left: Note that MAML is able to estimate parts of the curve where there are no datapoints, indicating that the model has learned about the periodic structure of sine waves. Right: Fine-tuning of a model pre-trained on the same distribution of tasks without MAML, with a tuned step size. Due to the often contradictory outputs on the pre-training tasks, this model is unable to recover a suitable representation and fails to extrapolate from the small number of test-time samples.]]<br />
<br />
=== Classification ===<br />
<br />
For classification evaluation, Omniglot and MiniImagenet datasets are used. The Omniglot dataset consists of 20 instances of 1623 characters from 50 different alphabets. <br />
<br />
The experiment involves fast learning of N-way classification with 1 or 5 shots. The problem of N-way classification is set up as follows: select N unseen classes, provide the model with K different instances of each of the N classes, and evaluate the model’s ability to classify new instances within the N classes. For Omniglot, 1200 characters are selected randomly for training, irrespective of alphabet, and use the remaining for testing. The Omniglot dataset is augmented with rotations by multiples of 90 degrees.<br />
<br />
The model follows the same architecture as the embedding function that has 4 modules with a 3-by-3 convolutions and 64 filters, followed by batch normalization, a ReLU nonlinearity, and 2-by-2 max-pooling. The Omniglot images are downsampled to 28-by-28, so the dimensionality of the last hidden layer is 64. The last layer is fed into a softmax. For Omniglot, strided convolutions is used instead of max-pooling. For MiniImagenet, 32 filters per layer are used to reduce overfitting. In order to also provide a fair comparison against memory-augmented neural networks [3] and to test the flexibility of MAML, the results for a non-convolutional network are also provided. <br />
<br />
[[File:ershad_results2.png|500px|center|thumb|Table 1: Few-shot classification on held-out Omniglot characters (top) and the MiniImagenet test set (bottom). MAML achieves results that are comparable to or outperform state-of-the-art convolutional and recurrent models. Siamese nets, matching nets, and the memory module approaches are all specific to classification, and are not directly applicable to regression or RL scenarios. The $\pm$ shows 95% confidence intervals over tasks. ]]<br />
<br />
=== Reinforcement Learning ===<br />
Several simulated continuous control environments are used for RL evaluation. In all of the domain, the MAML model is a neural network policy with two hidden layers of size 100, and ReLU activations. The gradient updates are computed using vanilla policy gradient and trust-region policy optimization (TRPO) is used as the meta-optimizer.<br />
<br />
In order to avoid computing third derivatives, finite differences are computed to <br />
compute the Hessian-vector products for TRPO. For both learning and meta-learning updates, we use the standard linear feature baseline proposed by [4], which is fitted separately at each iteration for each sampled task in the batch. <br />
<br />
Three baseline models for the comparison are: <br />
(a) pretraining one policy on all of the tasks and then fine-tuning<br />
(b) training a policy from randomly initialized weights<br />
(c) an oracle policy which receives the parameters of the task as input, which for the tasks below corresponds to a goal position, goal direction, or goal velocity for the agent. <br />
<br />
The baseline models of (a) and (b) are fine-tuned with gradient descent with a manually tuned step size.<br />
<br />
2D Navigation: In the first meta-RL experiment, the authors study a set of tasks where a point agent must move to different goal positions in 2D, randomly chosen for each task within a unit square. The observation is the current 2D position, and actions correspond to velocity commands clipped to be in the range [-0.1; 0.1]. The reward is the negative squared distance to the goal, and episodes terminate when the agent is within 0:01 of the goal or at the horizon ofH = 100. The policy was trained with MAML <br />
to maximize performance after 1 policy gradient update using 20 trajectories. They compare adaptation to a new task with up to 4 gradient updates, each with 40 samples. Results are shown in Fig. 3.<br />
<br />
[[File:ershad_results3.png|500px|center|thumb|Figure 3: Top: quantitative results from 2D navigation task, Bottom: qualitative comparison between model learned with MAML and with fine-tuning from a pretrained network ]]<br />
<br />
Locomotion. To study how well MAML can scale to more complex deep RL problems, we also study adaptation on high-dimensional locomotion tasks with the MuJoCo simulator [5]. The tasks require two simulated robots – a planar cheetah and a 3D quadruped (the “ant”) – to run in a particular direction or at a particular velocity. In the goal velocity experiments, the reward is the negative absolute value between the current velocity of the agent and a goal, which is chosen uniformly at random between 0 and 2 for the cheetah and between 0 and 3 for the ant. In the goal direction experiments, the reward is the magnitude of the velocity in either the forward or backward direction, chosen at random for each task in p(T ). The horizon is H = 200, with 20 rollouts per gradient step for all problems except the ant forward/backward task, which used 40 rollouts per step. The results in Figure 5 show that MAML learns a model that can quickly adapt its velocity and direction with even <br />
just a single gradient update, and continues to improve with more gradient steps. The results also show that, on these challenging tasks, the MAML initialization substantially outperforms random initialization and pretraining.<br />
[[File:ershad_results4.png|500px|center|thumb|Figure 4: Reinforcement learning results for the half-cheetah and ant locomotion tasks, with the tasks shown on the far right. ]]<br />
<br />
A conceptual method to achieve fast adaptation in language modeling tasks ( not been experimented on by the authors) would be to explore methods of attaching an Attention Kernel which results in a simple and differentiable loss. It has been implemented in One-Shot Language Modeling along with state-of-the-art improvements in one-shot learning on Imagenet and Omniglot [7].<br />
<br />
='''Conclusion'''=<br />
<br />
The paper introduced a meta-learning method based on learning easily adaptable model parameters through gradient descent. The approach has a number of benefits. It is simple and does not introduce any learned parameters for meta-learning. It can be combined with any model representation that is amenable to gradient-based training, and any differentiable objective, including classification, regression, and reinforcement learning. Lastly, since the method merely produces a weight initialization, adaptation can be performed with any amount of data and any number of gradient steps, though it demonstrates state-of-the-art results on classification with only one or five examples per class. The authors also show that the method can adapt an RL agent using policy gradients and a very modest amount of experience.<br />
<br />
='''Critique'''=<br />
From my opinion, the Model-Agnostic Meta-Learning looks like a simplified curriculum learning. It treats all tasks the same over the whole training history, and does not consider the difficulty of the tasks and the adaption of neural network to the task. Curriculum learning would be a good idea to speed up the training.<br />
<br />
<br />
='''References'''=<br />
# Schmidhuber, J¨urgen. Learning to control fast-weight memories: An alternative to dynamic recurrent networks. Neural Computation, 1992.<br />
# Lake, Brenden M, Salakhutdinov, Ruslan, Gross, Jason, and Tenenbaum, Joshua B. One shot learning of simple visual concepts. In Conference of the Cognitive Science Society (CogSci), 2011.<br />
# Santoro, Adam, Bartunov, Sergey, Botvinick, Matthew, Wierstra, Daan, and Lillicrap, Timothy. Meta-learning with memory-augmented neural networks. In International Conference on Machine Learning (ICML), 2016.<br />
# Duan, Yan, Chen, Xi, Houthooft, Rein, Schulman, John, and Abbeel, Pieter. Benchmarking deep reinforcement learning for continuous control. In International Conference on Machine Learning (ICML), 2016.<br />
# Todorov, Emanuel, Erez, Tom, and Tassa, Yuval. Mujoco: A physics engine for model-based control. In International Conference on Intelligent Robots and Systems (IROS), 2012.<br />
# Videos the learned policies can be found in https://sites.google.com/view/maml.<br />
# Oriol Vinyals, Charles Blundell, Timothy Lillicrap, Koray Kavukcuoglu, Daan Wierstra. "Matching Networks for One Shot Learning". arXiv:1606.04080 [cs.LG]<br />
<br />
Implementation Example: https://github.com/cbfinn/maml</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Exploration_via_Bootstrapped_DQN&diff=31111Deep Exploration via Bootstrapped DQN2017-11-22T18:48:59Z<p>Nkadiyar: /* Q Learning and Deep Q Networks [5] */</p>
<hr />
<div>== Details ==<br />
<br />
'''Title''': Deep Exploration via Bootstrapped DQN<br />
<br />
'''Authors''': Ian Osband {1,2}, Charles Blundell {2}, Alexander Pritzel {2}, Benjamin Van Roy {1}<br />
<br />
'''Organisations''':<br />
# Stanford University<br />
# Google Deepmind<br />
<br />
'''Conference''': NIPS 2016<br />
<br />
'''URL''': [https://papers.nips.cc/paper/6501-deep-exploration-via-bootstrapped-dqn papers.nips.cc]<br />
<br />
'''Online code sources'''<br />
* [https://github.com/iassael/torch-bootstrapped-dqn github.com/iassael/torch-bootstrapped-dqn]<br />
<br />
This summary contains background knowledge from Section 2-7 (except Section 5). Feel free to skip if you already know.<br />
<br />
== Intro to Reinforcement Learning ==<br />
<br />
In reinforcement learning, an agent interacts with an environment with the goal to maximize its long-term reward. A common application of reinforcement learning is to the [https://en.wikipedia.org/wiki/Multi-armed_bandit multi armed bandit problem]. In a multi armed bandit problem, there is a gambler and there are $n$ slot machines, and the gambler can choose to play any specific slot machine at any time. All the slot machines have their own probability distributions by which they churn out rewards, but this is unknown to the gambler. So the question is, how can the gambler learn the strategy to get the maximum long term reward?<br />
<br />
There are two things the gambler can do at any instance: either he can try a new slot machine, or he can play the slot machine he has tried before (and he knows he will get some reward). However, even though trying a new slot machine feels like it would bring less reward to the gambler, it is possible that the gambler finds out a new slot machine that gives a better reward than the current best slot machine. This is the dilemma of '''exploration vs exploitation'''. Trying out a new slot machine is '''exploration''', while redoing the best move so far is '''exploiting''' the currently understood perception of the reward.<br />
<br />
[[File:multiarmedbandit.jpg|thumb|Source: [https://blogs.mathworks.com/images/loren/2016/multiarmedbandit.jpg blogs.mathworks.com]]]<br />
<br />
There are many strategies to approach this '''exploration-exploitation dilemma'''. Some [https://web.stanford.edu/class/msande338/lec9.pdf common strategies] for optimizing in an exploration-exploitation setting are Random Walk, Curiosity-Driven Exploration, and Thompson Sampling. A lot of these approaches are provably efficient, but assume that the state space is not very large. For instance, the approach called Curiosity-Driven Exploration aims to take actions that lead to immediate additional information. This requires the model to search “every possible cell in the grid” which is not desirable if state space is very large. Strategies for large state spaces often just either ignore exploration, or do something naive like $\epsilon$-greedy, where you exploit with $1-\epsilon$ probability and explore "randomly" in rest of the cases. The general idea to tackle large or continuous state spaces is by value function approximation. An empirically tested strategy is Value Function Approximation using Fourier Basis [16]. It has also proven to perform well compared to radial basis functions and the polynomial basis, which are the two most popular fixed bases for linear value function approximation. <br />
<br />
This paper tries to use a Thompson sampling like approach to make decisions.<br />
<br />
== Thompson Sampling<sup>[[#References|[1]]]</sup> ==<br />
<br />
In Thompson sampling, our goal is to reach a belief that resembles the truth. Let's consider a case of coin tosses (2-armed bandit). Suppose we want to be able to reach a satisfactory pdf for $\mathbb{P}_h$ (heads). Assuming that this is a Bernoulli bandit problem, i.e. the rewards are $0$ or $1$, we can start off with $\mathbb{P}_h^{(0)}=\beta(1,1)$. The $\beta(x,y)$ distribution is a very good choice for a possible pdf because it works well for Bernoulli rewards. Further $\beta(1,1)$ is the uniform distribution $\mathbb{N}(0,1)$.<br />
<br />
Now, at every iteration $t$, we observe the reward $R^{(t)}$ and try to make our belief close to the truth by doing a Bayesian computation. Assuming $p$ is the probability of getting a heads,<br />
<br />
$$<br />
\begin{align*}<br />
\mathbb{P}(R|D) &\propto \mathbb{P}(D|R) \cdot \mathbb{P}(R) \\<br />
\mathbb{P}_h^{(t+1)}&\propto \mbox{likelihood}\cdot\mbox{prior} \\<br />
&\propto p^{R^{(t)}}(1-p)^{R^{(t)}} \cdot \mathbb{P}_h^{(t)} \\<br />
&\propto p^{R^{(t)}}(1-p)^{R^{(t)}} \cdot \beta(x_t, y_t) \\<br />
&\propto p^{R^{(t)}}(1-p)^{R^{(t)}} \cdot p^{x_t-1}(1-p)^{y_t-1} \\<br />
&\propto p^{x_t+R^{(t)}-1}(1-p)^{y_t+R^{(t)}-1} \\<br />
&\propto \beta(x_t+R^{(t)}, y_t+R^{(t)})<br />
\end{align*}<br />
$$<br />
<br />
[[File:thompson sampling coin example.png|thumb||||600px|Source: [https://www.quora.com/What-is-Thompson-sampling-in-laymans-terms Quora]]]<br />
<br />
This means that with successive sampling, our belief can become better at approximating the truth. There are similar update rules if we use a non Bernoulli setting, say, Gaussian. In the Gaussian case, we start with $\mathbb{P}_h^{(0)}=\mathbb{N}(0,1)$ and given that $\mathbb{P}_h^{(t)}\propto\mathbb{N}(\mu, \sigma)$ it is possible to show that the update rule looks like<br />
<br />
$$<br />
\mathbb{P}_h^{(t+1)} \propto \mathbb{N}\bigg(\frac{t\mu+R^{(t)}}{t+1},\frac{\sigma}{\sigma+1}\bigg)<br />
$$<br />
<br />
=== How can we use this in reinforcement learning? ===<br />
<br />
We can use this idea to decide when to explore and when to exploit. We start with an initial belief, choose an action, observe the reward and based on the kind of reward, we update our belief about what action to choose next.<br />
<br />
== Bootstrapping <sup>[[#References|[2,3]]]</sup> ==<br />
<br />
This idea may be unfamiliar to some people, so I thought it would be a good idea to include this. In statistics, bootstrapping is a method to generate new samples from a given sample. Suppose that we have a given population, and we want to study a property $\theta$ of the population. So, we just find $n$ sample points (sample $\{D_i\}_{i=1}^n$), calculate the estimator of the property, $\hat{\theta}$, for these $n$ points, and make our inference. <br />
<br />
If we later wish to find some property related to the estimator $\hat{\theta}$ itself, e.g. we want a bound of $\hat{\theta}$ such that $\delta_1 \leq \hat{\theta} \leq \delta_2$ with a confidence of $c=95%$, then we can use bootstrapping for this.<br />
<br />
Using bootstrapping, we can create a new sample $\{D'_i\}_{i=1}^{n'}$ by '''randomly sampling $n'$ times from $D$, with replacement'''. So, if $D=\{1,2,3,4\}$, a $D'$ of size $n'=10$ could be $\{1,4,4,3,2,2,2,1,3,4\}$. We do this a sufficient $k$ number of times, calculate $\hat{\theta}$ each time, and thus get a distribution $\{\hat{\theta}_i\}_{i=1}^k$. Now, we can choose the $100\cdot c$<sup>th</sup> and $100\cdot(1-c)$<sup>th</sup> percentile of this distribution, (let them be $\hat{\theta}_\alpha$ and $\hat{\theta}_\beta$ respectively) and say<br />
<br />
$$\hat{\theta}_\alpha \leq \hat{\theta} \leq \hat{\theta}_\beta, \mbox{with confidence }c$$<br />
<br />
== Why choose bootstrap and not dropout? ==<br />
<br />
There is previous work<sup>[[#References|[4]]]</sup> that establishes dropout as a good way to train NNs on a posterior such that the trained NN works like a function approximator that is close to the actual posterior. But, there are several problems with the predictions of this trained NN. The figures below are from the appendix of this paper. The left image is the NN trained by the authors of this paper on a sample noisy distribution and the right image is from the accompanying web demo from [[#References|[4]]], where the authors of [[#References|[4]]] show that their NN converges around the mean with a good confidence.<br />
<br />
[[File:dropout_results.png|thumb||center||700px|Source: this paper's appendix]]<br />
<br />
According to the authors of this paper,<br />
# Even though [[#References|[4]]] says that dropout converges arond the mean, their experiment actually behaves weirdly around a reasonable point like $x=0.75$. They think that this happens because dropout only affects the region local to the original data.<br />
# Samples from the NN trained on the original data do not look like a reasonable posterior (very spiky).<br />
# The trained NN collapses to zero uncertainty at the data points from the original data.<br />
<br />
== Q Learning and Deep Q Networks <sup>[[#References|[5]]]</sup> ==<br />
<br />
At any point of time, our rewards dictate what our actions should be. Also, in general, we want good long term rewards. For example, if we are playing a first person shooter game, it is a good idea to go out of cover to kill an enemy, even if some health is lost. Similarly, in reinforcement learning, we want to maximize our long term reward. So if at each time $t$, the reward is $r_t$, then a naive way is to say we want to maximise<br />
<br />
$$<br />
R_t = \sum_{i=0}^{\infty}r_t<br />
$$<br />
<br />
But, this reward is unbounded. So technically it could tend to $\infty$ in a lot of the cases. This is why we use a '''discounted reward'''.<br />
<br />
$$<br />
R_t = \sum_{i=0}^{\infty}\gamma^t r_t<br />
$$<br />
<br />
Here, we take $0\leq \gamma \lt 1$. If it is equal to one, the agent values future reward just as much as current reward. Conversely, a value of zero will cause the agent to only value immediate rewards, which only works with very detailed reward functions. So, what this means is that we value our current reward the most ($r_0$ has a coefficient of $1$), but we also consider the future possible rewards. So if we had two choices: get $+4$ now and $0$ at all other timesteps, or get $-2$ now and $+2$ after $3$ timesteps for $20$ timesteps, we choose the latter ($\gamma=0.9$). This is because $(+4) < (-2)+0.9^3(2+0.9\cdot2+\cdots+0.9^{19}\cdot2)$.<br />
<br />
<br />
A '''policy''' $\pi: \mathbb{S} \rightarrow \mathbb{A}$ is just a function that tells us what action to take in a given state $s\in \mathbb{S}$. Our goal is to find the best policy $\pi^*$ that maximises the reward from a given state $s$. So, a '''value function''' is defined from $s$ (which the agent is in, at timestep $t$) and following the policy $\pi$ as $V^\pi(s) = \mathbb{E}[R_t]$. The optimal value function is then simply<br />
<br />
$$<br />
V^*(s)=\displaystyle\max_{\pi}V^\pi(s)<br />
$$<br />
<br />
For convenience however, it is better to work with the '''Q function''' $Q: \mathbb{S}\times\mathbb{A} \rightarrow \mathbb{R}$. $Q$ is defined similarly as $V$. It is the expected return after taking an action $a$ in the given state $s$. So, $Q^\pi(s,a)=\mathbb{E}[R_t|s,a]$. The optimal $Q$ function is<br />
<br />
$$<br />
Q^*(s,a)=\displaystyle\max_{\pi}Q^\pi(s,a)<br />
$$<br />
<br />
Suppose that we know $Q^*$. Then, if we know that we are supposed to start at $s$ and take an action $a$ right now, what is the best course of action from the next time step? We just choose the optimal action $a'$ at the next state $s'$ that we reach. The optimal action $a'$ at state $s'$ is simply the argument $a_x$ that maximises our $Q^*(s',\cdot)$.<br />
<br />
$$<br />
a'=\displaystyle\arg\max_{a_x} Q^*(s',a_x)<br />
$$<br />
<br />
So, our best expected reward from $s$ taking action $a$ is $\mathbb{E}[r_t+\gamma\mathbb{E}[R_{t+1}]]$. This is known as the '''Bellman equation''' in optimal control problem (By the way, its continuous form is called '''Hamilton-Jacobi-Bellman equation''' or HJB equation, which is a very important partial differential equation):<br />
<br />
$$<br />
Q^*(s,a)=\mathbb{E}[r_t+\gamma \displaystyle\max_{a_x} Q^*(s',a_x)]<br />
$$<br />
<br />
In Q learning, we use a deep neural network with weights $\theta$ as a function approximator for $Q^*$, since Bellman equation is indeed a non-linear PDE and very difficult to solve numerically. The '''naive way''' to do this is to design a deep neural network that takes as input the state $s$ and action $a$, and produces an approximation to $Q^*$. <br />
<br />
* Suppose our neural net weights are $\theta_i$ at iteration $i$.<br />
* We want to train our neural net on the case when we are at $s$, take action $a$, get reward $r$, and reach $s'$.<br />
* To find out what action is best from $s'$, i.e. $a'$, we have to simulate all actions from $s'$. We can do this after we complete this iteration, then run $s',a_x$ for all $a_x\in\mathbb{A}$. But, we don't know how to complete this iteration without knowing this $a'$. So, another way is to simulate all actions from $s'$ using last known set of weights $\theta_{i-1}$. We just simulate state $s'$, action $a_x$ for all $a_x\in\mathbb{A}$ from the previous state and get $Q^*(s',a_x;\theta_{i-1})$. ('''Note''' that some papers do not use the set of weights from the previous iteration $\theta_{i-1}$. Instead they fix the weights for finding the best action for every $\tau$ steps to $\theta^-$, and do $Q^*(s',a_x;\theta^-)$ for $a_x\in\mathbb{A}$ and use this for the target value.)<br />
* Now we can compute our loss function using the Bellman equation, and backpropagate.<br />
$$<br />
\mbox{loss}=\mbox{target}-\mbox{prediction}=(r+\displaystyle\max_{a_x}Q^*(s',a_x;\theta_{i-1}))-Q^*(s,a;\theta_i)<br />
$$<br />
<br />
The '''problem''' with this approach is that at every iteration $i$, we have to do $|\mathbb{A}|$ forward passes on the previous set of weights $\theta_{i-1}$ to find out the best action $a'$ at $s'$. This becomes infeasible quickly with more possible actions.<br />
<br />
Authors of [[#References|[5]]] therefore use another kind of architecture. This architecture takes as input the state $s$, and computes the values $Q^*(s,a_x)$ for $a_x\in\mathbb{A}$. So there are $|\mathbb{A}|$ outputs. This basically parallelizes the forward passes so that $r+\displaystyle\max_{a_x}Q^*(s',a_x;\theta_{i-1})$ can be done with just a single pass through the outputs.<br />
<br />
<br />
[[File:DQN_arch.png|thumb||||600px|Source: [https://leonardoaraujosantos.gitbooks.io/artificial-inteligence/content/image_folder_7/DQNBreakoutBlocks.png leonardoaraujosantos.gitbooks.io]]]<br />
<br />
'''Note:''' When I say state $s$ as an input, I mean some representation of $s$. Since the environment is a partially observable MDP, it is hard to know $s$. So, we can for example, apply a CNN on the frames and get an idea of what the current state is. We pass this output to the input of the DNN (DNN is the fully connected layer for the CNN then).<br />
<br />
=== Experience Replay ===<br />
<br />
Authors of this paper borrow the concept of experience replay from [[#References|[5,6]]]. In experience replay, we do training in episodes. In each episode, we play and store consecutive $(s,a,r,s')$ tuples in the experience replay buffer. Then after the play, we choose random samples from this buffer and do our training.<br />
<br />
<br />
Advantages of experience replay over simple online Q learning<sup>[[#References|[5]]]</sup>:<br />
* '''Better data efficiency''': It is better to use one transition many times to learn again and again, rather than just learn once from it.<br />
* Learning from consecutive samples is difficult because of correlated data. Experience replay breaks this correlation.<br />
* Online learning means the input is decided by the previous action. So, if the maximising action is to go left in some game, next inputs would be about what happens when we go left. This can cause the optimiser to get stuck in a feedback loop, or even diverge, as [[#Reference|[7]]] points out.<br />
<br />
== Double Q Learning ==<br />
<br />
=== Problem with Q Learning<sup>[[#References|[8]]]</sup> ===<br />
<br />
For a simple neural network, each update tries to shift the current $Q^*$ estimate to a new value:<br />
<br />
$$<br />
Q^*(s,a) \leftarrow r+\gamma\displaystyle\max_{a_x}Q^*(s',a_x)<br />
$$<br />
<br />
Suppose the neural net has some inherent noise $\epsilon$. So, the neural net actually stores a value $\mathbb{Q}^*$ given by<br />
<br />
$$<br />
\mathbb{Q}^* = Q^*+\epsilon<br />
$$<br />
<br />
Even if $\epsilon$ has zero mean in the beginning, using the $\max$ operator at the update steps will start propagating $\gamma\cdot\max \mathbb{Q}^*$. This leads to a non zero mean subsequently. The problem is that "max causes overestimation because it does not preserve the zero-mean property of the errors of its operands." ([[#References|[8]]]) Thus, Q learning is more likely to choose overoptimistic values.<br />
<br />
=== How does Double Q Learning work? <sup>[[#References|[9]]]</sup> ===<br />
<br />
The problem can be solved by using two sets of weights $\theta$ and $\Theta$. The $\mbox{target}$ can be broken up as<br />
<br />
$$<br />
\mbox{target} = r+\displaystyle\max_{a_x}Q^*(s',a_x;\theta) = r+Q^*(s',\displaystyle\arg\max_{a_x}Q^*(s',a_x;\theta);\theta) = r+Q^*(s',a';\theta)<br />
$$<br />
<br />
Using double Q learning, we '''select''' the best action using current weights $\theta$ and '''evaluate''' the $Q^*$ value to decide the target value using $\Theta$.<br />
<br />
$$<br />
\mbox{target} = r+Q^*(s',\displaystyle\arg\max_{a_x}Q^*(s',a_x;\theta);\Theta) = r+Q^*(s',a';\Theta)<br />
$$<br />
<br />
This makes the evaluation fairer.<br />
<br />
=== Double Deep Q Learning ===<br />
<br />
[[#References|[9]]] further talks about how to use this for deep learning without much additional overhead. The suggestion is to use $\theta^-$ as $\Theta$.<br />
<br />
$$<br />
\mbox{target} = r+Q^*(s',\displaystyle\arg\max_{a_x}Q^*(s',a_x;\theta);\theta^-) = r+Q^*(s',a';\theta^-)<br />
$$<br />
<br />
== Bootstrapped DQN ==<br />
<br />
The authors propose an architecture that has a shared network and $K$ bootstrap heads. So, suppose our experience buffer $E$ has $n$ data points, where each datapoint is a $(s,a,r,s')$ tuple. Each bootstrap head trains on a different buffer $E_i$, where each $E_i$ has been constructed by sampling $n$ data points from the original experience buffer $E$ with replacement ('''bootstrap method''').<br />
<br />
<br />
Because each of the heads train on a different buffer, they model a different $Q^*$ function (say $Q^*_k$). Now, for each episode, we first choose a specific $Q^*_k=Q^*_s$. This $Q^*_s$ helps us create the experience buffer for the episode. From any state $s_t$, we populate the experience buffer by choosing the next action $a_t$ that maximises $Q^*_s$. (similar to '''Thompson Sampling''')<br />
<br />
$$<br />
a_t = \displaystyle\arg\max_a Q^*_s(s_t,a_t)<br />
$$<br />
<br />
Also, along with $s_t,a_t,r_t,s_{t+1}$, they push a bootstrap mask $m_t$. This mask is basically is a binary vector of size $K$, and it tells which $Q_k$ should be affected by this datapoint, if it is chosen as a training point. So, for example, if $K=5$ and there is a experience tuple $(s_t,a_t,r_t,s_{t+1},m_t)$ where $m_t=(0,1,1,0,1)$, then $(s_t,a_t,r_t,s_{t+1})$ should only affect $Q_2,Q_3$ and $Q_5$.<br />
<br />
<br />
So, at each iteration, we just choose few points from this buffer and train the respective $Q_{(\cdot)}$ based on the bootstrap masks.<br />
<br />
=== How to generate masks? ===<br />
<br />
Masks are created by sampling from the '''masking distribution'''. Now, there are many ways to choose this masking distribution:<br />
<br />
* If for each datapoint $D_i$ ($i=1$ to $n$), we mask from $\mbox{Bernoulli}(0.5)$, this will roughly allow us to have half the points from the original buffer. To get to size $n$, we duplicate these points by doubling the weights for each datapoint. This essentially gives us a '''double or nothing''' bootstrap<sup>[[#References|[10]]]</sup>.<br />
* If the mask is $(1, 1 \cdots 1)$, then this becomes an '''ensemble learning''' method.<br />
* $m_t~\mbox{Poi}(1)$ (poisson distribution)<br />
* $m_t[k]~\mbox{Exp}(1)$ (exponential distribution)<br />
<br />
For this paper's results, the authors used a $\mbox{Bernoulli}(p)$ distribution.<br />
<br />
== Related Work ==<br />
<br />
The authors mention the method described in [[#References|[11]]]. The authors of [[#References|[11]]] talk about the principle of "optimism in the face of uncertainty" and modify the reward function to encourage state-action pairs that have not been seen often:<br />
<br />
$$<br />
R(s,a) \leftarrow R(s,a)+\beta\cdot\mbox{novelty}(s,a)<br />
$$<br />
<br />
According to the authors, [[#References|[11]]]'s DQN algorithm relies on a lot of hand tuning and is only good for non stochastic problems. The authors further compare their results to [[#References|[11]]]'s results on Atari.<br />
<br />
<br />
The authors also mention an existing algorithm PSRL<sup>[[#References|[12,13]]]</sup>, or posterior sampling based RL. However, this algorithm requires a solved MDP, which is not feasible for large systems. Bootstrapped DQN approximates this idea by sampling from approximate $Q^*$ functions.<br />
<br />
<br />
Further, the authors mention that the work in [[#References|[12,13]]] has been followed by RLSVI<sup>[[#Reference|[14]]]</sup> which solves the problem for linear cases.<br />
<br />
== Deep Exploration: Why is Bootstrapped DQN so good at it? ==<br />
<br />
The authors consider a simple example to demonstrate the effectiveness of bootstrapped DQN at deep exploration.<br />
<br />
[[File:deep_exploration_example.png|thumb||center||700px|Source: this paper, section 5.1]]<br />
<br />
<br />
<br />
In this example, the agent starts at $s_2$. There are $N$ steps, and $N+9$ timesteps to generate the experience buffer. The agent is said to have learned the optimal policy if it achieves the best possible reward of $10$ (go to the rightmost state in $N-1$ timesteps, then stay there for $10$ timesteps), for at least $100$ such episodes. The results they got:<br />
<br />
[[File:deep_exploration_results.png|thumb||center||700px|Source: this paper, section 5.1]]<br />
<br />
<br />
<br />
The blue dots indicate when the agent learnt the optimal policy. If this took more than $2000$ episodes, they indicate it with a red dot. Thompson DQN is DQN with posterior sampling at every timestep. Ensemble DQN is same as bootstrapped DQN except that the mask is all $(1,1 \cdots 1)$. It is evident from the graphs that bootstrapped DQN can achieve deep exploration better than these two methods, and DQN.<br />
<br />
=== But why is it better? ===<br />
<br />
The authors say that this is because bootstrapped DQN constructs different approximations to the posterior $Q^*$ with the same initial data. This diversity of approximations is because of random initalization of weights for the $Q^*_k$ heads. This means that these heads start out trying random actions (because of diverse random initial $Q^*_k$), but when some head finds a good state and generalises to it, some (but not all) of the heads will learn from it, because of the bootstrapping. Eventually other heads will either find other good states, or end up learning the best good states found by the other heads.<br />
<br />
<br />
So, the architecture explores well and once a head achieves the optimal policy, eventually, all heads achieve the policy.<br />
<br />
== Results ==<br />
<br />
The authors test their architecture on 49 Atari games. They mention that there has been recent work to improve the performance of DDQNs, but those are tweaks whose intentions are orthogonal to this paper's idea. So, they don't compare their results with them.<br />
<br />
=== Scale: What values of $K$, $p$ are best? ===<br />
<br />
[[File:scale_k_p.png|thumb||center||800px|Source: this paper, section 6.1]]<br />
<br />
Recall that $K$ is the number of bootstrap heads and $p$ is the parameter for the masking distribution (Bernoulli). The authors say that around $K=10$, the performance reaches close to the peak, so it should be good.<br />
<br />
<br />
$p$ also represents the amount of data sharing. This is because lesser $p$ means there is lesser chance (due to the Bernoulli distribution) that the corresponding datapoint is taken into the bootstrapped dataset $D_i$. So, lesser $p$ means more identical datapoints, hence more heads share their datapoints. However, the value of $p$ doesn't seem to affect the rewards achieved over time. The authors give the following reasons for it:<br />
<br />
* The heads start with random weights for $Q^*$, so the targets (which use $Q^*$) turn out to be different. So the update rules are different.<br />
* Atari is deterministic.<br />
* Because of the initial diversity, the heads will learn differently even if they predict the same action for the given state.<br />
<br />
$p=1$ is the value they use finally, because this reduces the no. of identical datapoints and reduces time.<br />
<br />
=== Performance on Atari ===<br />
<br />
In general, the results tell us that bootstrapped DQN achieves better results.<br />
<br />
[[File:atari_results_bootstrapped_dqn.png|thumb||center||800px|Source: this paper, section 6.2]]<br />
<br />
The authors plot the improvement they achieved with bootstrapped DQN with the games. They define '''improvement''' to be $x$ if bootstrapped DQN achieves a better result than DQN in $\frac{1}{x}$ frames.<br />
<br />
[[File:bdqn_improvement.png|thumb||center||1000px|Source: this paper, section 6.2]]<br />
<br />
<br />
The authors say that bootstrapped DQN doesn't work good on all Atari games. They point out that there are some challenging games, where exploration is key but bootstrapped DQN doesn't do good enough (but does better than DQN). Some of these games are Frostbite and Montezuma’s Revenge. They say that even better exploration may help, but also point out that there may be other problems like: network instability, reward clipping and temporally extended rewards.<br />
<br />
=== Improvement: Highest Score Reached & how fast is this high score reached? ===<br />
<br />
The authors plot the improvement graphs after 20m and 200m frames.<br />
<br />
[[File:cumulative_rewards_bdqn.png|thumb||center||700px|Source: this paper, section 6.3]]<br />
<br />
=== Visualisation of Results ===<br />
<br />
One of the authors' [https://www.youtube.com/playlist?list=PLdy8eRAW78uLDPNo1jRv8jdTx7aup1ujM youtube playlist] can be found online.<br />
<br />
<br />
The authors also point out that just purely using bootstrapped DQN as an exploitative strategy is pretty good by itself, better than vanilla DQN. This is because of the deep exploration capabilities of bootstrapped DQN, since it can use the best states it knows and also plan to try out states it doesn't have any information about. Even in the videos, it can be seen that the heads agree at all the crucial decisions, but stay diverse at other less important steps.<br />
<br />
== Critique ==<br />
<br />
It would be very interesting and a great addition to the experimental section of the paper, if the authors would have compared with asynchronous methods of exploration of the state space first introduced in [[#References|[15]]]. The authors unfortunately only compared their DQN with the original DQN and not all the other variations in the literature, and justified it by saying that their idea was "orthogonal" to these improvements.<br />
<br />
=== Different way to do exploration-exploitation? ===<br />
<br />
Instead of choosing the next action $a_t$ that maximises $Q^*_s$, they could have chosen different actions $a_i$ with probabilities<br />
<br />
$$<br />
\mathbb{P}(s_t,a_i) = \frac{Q^*_s(s_t,a_i)}{\displaystyle \sum_{i=1}^{|\mathbb{A}|} Q^*_s(s_t,a_i)}<br />
$$<br />
<br />
According to me, this is closer to Thompson Sampling.<br />
<br />
=== Why use Bernoulli? ===<br />
<br />
The choice of having a Bernoulli masking distribution eventually doesn't help them at all, since the algorithm does well because of the initial diversity. Maybe they can use some other masking distribution? However, the bootstrapping procedure is distribution-independent, the choice of masking distribution should not affect the long term performance of Bootstrapped DQN.<br />
<br />
=== Unanswered Questions & Miscellaneous ===<br />
* The Thompson DQN is not preferred because other randomized value functions can implement settings similar to Thompson sampling without the need for an intractable exact posterior update and also by working around the computational issue with Thompson Sampling: resampling every time step. Perhaps the authors could have explored Temporal Difference learning which is an attempt at combining Dynamic Programming and Monte Carlo methods.<br />
* The actual algorithm is hidden in the appendix. It could have been helpful if it were in the main paper.<br />
<br />
=== Improvement on Hierarchical Tasks ===<br />
In this paper, it is interesting that such bootstrap exploration principle actually improves the performance over hierarchical tasks such as Montezuma Revenge. It would be better if the authors could illustrate further about the influence of exploration in a sparse-reward hierarchical task.<br />
<br />
== References ==<br />
<br />
# [https://bandits.wikischolars.columbia.edu/file/view/Lecture+4.pdf Learning and optimization for sequential decision making, Columbia University, Lec 4]<br />
# [https://www.thoughtco.com/what-is-bootstrapping-in-statistics-3126172 Thoughtco, What is bootstrapping in statistics?]<br />
# [https://ocw.mit.edu/courses/mathematics/18-05-introduction-to-probability-and-statistics-spring-2014/readings/MIT18_05S14_Reading24.pdf Bootstrap confidence intervals, Class 24, 18.05, MIT Open Courseware]<br />
# [https://arxiv.org/abs/1506.02142 Yarin Gal and Zoubin Ghahramani. Dropout as a bayesian approximation: Representing model uncertainty in deep learning. arXiv preprint arXiv:1506.02142, 2015.]<br />
# [https://www.cs.toronto.edu/~vmnih/docs/dqn.pdf Mnih et al., Playing Atari with Deep Reinforcement Learning, 2015]<br />
# Long-Ji Lin. Reinforcement learning for robots using neural networks. Technical report, DTIC Document, 1993.<br />
# John N Tsitsiklis and Benjamin Van Roy. An analysis of temporal-difference learning with function approximation. Automatic Control, IEEE Transactions on, 42(5):674–690, 1997.<br />
# S. Thrun and A. Schwartz. Issues in using function approximation for reinforcement learning, 1993.<br />
# [https://arxiv.org/pdf/1509.06461.pdf Hado Van Hasselt, Arthur Guez, and David Silver. Deep reinforcement learning with double q-learning, 2015.]<br />
# [https://pdfs.semanticscholar.org/d623/c2cbf100d6963ba7dafe55158890d43c78b6.pdf Dean Eckles and Maurits Kaptein, Thompson Sampling with the Online Bootstrap, 2014, Pg 3]<br />
# [https://arxiv.org/abs/1507.00814 Bradly C. Stadie, Sergey Levine, Pieter Abbeel, Incentivizing Exploration In Reinforcement Learning With Deep Predictive Models, 2015.]<br />
# Ian Osband, Daniel Russo, and Benjamin Van Roy. (More) efficient reinforcement learning via posterior sampling, NIPS 2013.<br />
# Ian Osband and Benjamin Van Roy. Model-based reinforcement learning and the eluder dimension, NIPS 2014.<br />
# [https://arxiv.org/abs/1402.0635 Ian Osband, Benjamin Van Roy, Zheng Wen, Generalization and Exploration via Randomized Value Functions, 2014.]<br />
# Mnih, Volodymyr, et al. "Asynchronous methods for deep reinforcement learning." International Conference on Machine Learning. 2016.<br />
# George Konidaris, Sarah Osentoski, and Philip Thomas. 2011. Value function approximation in reinforcement learning using the fourier basis. In Proceedings of the Twenty-Fifth AAAI Conference on Artificial Intelligence (AAAI'11). AAAI Press 380-385.<br />
<br />
<br />
Other helpful links (unsorted):<br />
* [http://pemami4911.github.io/paper-summaries/deep-rl/2016/08/16/Deep-exploration.html pemami4911.github.io]<br />
* [http://www.stat.yale.edu/~pollard/Courses/241.fall97/Poisson.pdf Poisson Approximations]<br />
<br />
== Appendix ==<br />
<br />
=== Algorithm for Bootstrapped DQN ===<br />
The appendix lists the following algorithm. Periodically, the replay buffer is played back to update value function network Q.<br />
<br />
[[File:alg1.PNG|thumb||left||700px|Source: this paper's appendix]]</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Unsupervised_Domain_Adaptation_with_Residual_Transfer_Networks&diff=31012Unsupervised Domain Adaptation with Residual Transfer Networks2017-11-21T17:30:32Z<p>Nkadiyar: </p>
<hr />
<div>== Introduction ==<br />
'''Domain Adaptation''' [https://en.wikipedia.org/wiki/Domain_adaptation]is a problem in machine learning which involves taking a model which has been trained on a source domain, and applying this to a different (but related) target domain. '''Unsupervised domain adaptation''' refers to the situation in which the source data is labelled, while the target data is (predominantly) unlabeled. The problem at hand is then finding ways to generalize the learning on the source domain to the target domain. In the age of deep networks this problem has become particularly salient due to the need for vast amounts of labeled training data, in order to reap the benefits of deep learning. Manual generation of labeled data is often prohibitive, and in absence of such data networks are rarely performant. The attempt to circumvent this drought of data typically necessitates the gathering of "off-the-shelf" data sets, which are tangentially related and contain labels, and then building models in these domains. The fundamental issue that unsupervised domain adaptation attempts to address is overcoming the inherent shift in distribution across the domains, without the ability to observe this shift directly.<br />
<br />
This paper proposes a method for unsupervised domain adaptation which relies on three key components: <br />
# A kernel-based penalty to ensure that the abstract representations generated by the networks hidden layers are similar between the source and the target data; <br />
# An entropy based penalty on the target classifier, which exploits the entropy minimization principle; and <br />
# A residual network structure is appended, which allows the source and target classifiers to differ by a (learned) residual function, thus relaxing the shared classifier assumption which is traditionally made.<br />
<br />
This method outperforms state-of-the-art techniques on common benchmark datasets, and is flexible enough to be applied in most feed-forward neural networks.<br />
<br />
[[File:Source-and-Target-Domain-Office-31-Backpack.png|thumb|right|The Office-31 Dataset Images for Backpack. Shows the variation in the source and target domains to motivate why these methods are important.]] <br />
=== Working Example (Office-31) === <br />
In order to assist in the understanding of the methods, it is helpful to have a tangible sense of the problem front of mind. The Domain Adaptation Project [https://people.eecs.berkeley.edu/~jhoffman/domainadapt/] provides data sets which are tailored to the problem of unsupervised domain adaptation. One of these data sets (which is later used in the experiments of this paper) has images which are labeled based on the Amazon product page for the various items. There are then corresponding pictures taken either by webcams or digital SLR cameras. The goal of unsupervised domain adaptation on this data set would be to take any of the three image sources as the source domain, and transfer a classifier to the other domain; see the example images to understand the differences.<br />
<br />
One can imagine that, while it is likely easy to scrape labeled images from Amazon, it is likely far more difficult to collect labeled images from webcam or DSLR pictures directly. The ultimate goal of this method would be to train a model to recognize a picture of a backpack taken with a webcam, based on images of backpacks scraped from Amazon (or similar tasks).<br />
<br />
== Related Work ==<br />
Broadly speaking, the problem of domain adaptation mitigates manual labeling of data in areas such as machine learning, computer vision, and natural language processing. The general goal of domain adaptation is to reduce the discrepancy in probability distributions between the source and target domains.<br />
<br />
Research into the use of Deep Neural Networks for the purpose of domain adaptation has suggested that, while networks learn abstract feature representations which can reduce the discrepancy across domains, it is not possible to wholly remove it [http://www.icml-2011.org/papers/342_icmlpaper.pdf], [https://arxiv.org/pdf/1412.3474.pdf]. Further work has been done to design networks which adapt traditional deep nets (typically CNNs) to specifically address the problems posed by domain adaptation, these methods all only address the issue of feature adaptation [https://arxiv.org/pdf/1502.02791.pdf], [https://arxiv.org/pdf/1409.7495.pdf], [https://people.eecs.berkeley.edu/~jhoffman/papers/Tzeng_ICCV2015.pdf]. That is, they all assume that the target and source classifiers are shared between domains. <br />
<br />
The authors drew particular motivation from He et al. [https://arxiv.org/abs/1512.03385] with the proposed structure of residual networks. Combining the insights from the ResNet architecture, in addition to previous work that had leveraged classifier adaptation (in the context where some target data is labeled) [http://citeseerx.ist.psu.edu/viewdoc/download?doi=10.1.1.130.8224&rep=rep1&type=pdf], [http://www.machinelearning.org/archive/icml2009/papers/445.pdf], [http://ieeexplore.ieee.org/document/5539870/] the authors develop their proposed network.<br />
<br />
== Residual Transfer Networks ==<br />
Generally, in an unsupervised domain adaptation problem, we are dealing with a set $\mathcal{D}_s$ (called the source domain) which is defined by $\{(x_i^s, y_i^s)\}_{i=1}^{n_s}$. That is the set of all labeled input-output pairs in our source data set. We denote the number of source elements by $n_s$. There is a corresponding set $\mathcal{D}_t = \{(x_i^t)\}_{i=1}^{n_t}$ (the target domain), consisting of unlabeled input values. There are $n_t$ such values. <br />
[[File:RTN-Structure.png|thumb|left|upright|The overarching structure of the RTN. Consists of an existing network, to which a bottleneck, MMD block, and residual block is appended.]]<br />
We can think of $\mathcal{D}_s$ as being sampled from some underlying distribution $p$, and $\mathcal{D}_t$ as being sampled from $q$. Generally we have that $p \neq q$, partially motivating the need for domain adaptation methods. <br />
<br />
We can consider the classifiers $f_s(\underline{x})$ and $f_t(\underline{x})$, for the source domain and target domain respectively. It is possible to learn $f_s$ based on the sample $\mathcal{D}_s$. Under the '''shared classifier assumption''' it would be the case that $f_s(\underline{x}) = f_t(\underline{x})$, and thus learning the source classifier is enough. This method relaxes this assumption, assuming that in general $f_s \neq f_t$, and attempting to learn both.<br />
<br />
The example network extends deep convolutional networks (in this case AlexNet [http://vision.stanford.edu/teaching/cs231b_spring1415/slides/alexnet_tugce_kyunghee.pdf]) to '''Residual Transfer Networks''', the mechanics of which are outlined below. Recall that, if $L(\cdot, \cdot)$ is taken to be the cross-entropy loss function, then the empirical error of a CNN on the source domain $\mathcal{D}_s$ is given by:<br />
<br />
<center><br />
<math display="block"><br />
\min_{f_s} \frac{1}{n_s} \sum_{i=1}^{n_s} L(f_s(x_i^s), y_i^s)<br />
</math> <br />
</center><br />
<br />
In a standard implementation, the CNN optimizes over the above loss. This will be the starting point for the RTN.<br />
<br />
=== Structural Overview ===<br />
The model proposed in this paper extends existing CNN's and alters the loss function that is optimized over. While each of these components is discussed in depth below, the overarching architecture involves four components:<br />
<br />
# An existing deep model. While this can be any model, in theory, the authors leverage AlexNet in practice.<br />
# A bottleneck layer, used to reduce the dimensionality of the learned abstract feature space, directly after the existing network.<br />
# An MMD block, with the expressed intention of feature adaptation.<br />
# A residual block, with the expressed intention of classifier adaptation. <br />
<br />
This structure is then optimized against a loss function which combines the standard cross-entropy penalty with MMD and target entropy penalties, yielding the proposed Residual Transfer Network (RTN) structure.<br />
<br />
=== Feature Adaptation ===<br />
Feature adaptation refers to the process in which the features which are learned to represent the source domain are made applicable to the target domain. Broadly speaking a CNN works to generate abstract feature representations of the distribution that the inputs are sampled from. It has been found that using these deep features can reduce, but not remove, cross-domain distribution discrepancy, hence the need for feature adaptation. It is important to note that CNN's transfer from general to specific features as the network gets deeper. In this light, the discrepancy between the feature representation of the source and the target will grow through a deeper convolutional net. As such a technique for forcing these distributions to be similar is needed.<br />
<br />
In particular the authors of this paper impose a bottleneck layer (call it $fc_b$) which is included after the final convolutional layer of AlexNet. This dense layer is connected to an additional dense layer $fc_c$, (which will serve as the target classification layer). They then compute the tensor product between the activations of the layers, performing "lossless multi-layer feature fusion". That is for the source domain they define $z_i^s \overset{\underset{\mathrm{def}}{}}{=} x_i^{s,fc_b}\otimes x_i^{s,fc_c}$ and for the target domain, $z_i^t \overset{\underset{\mathrm{def}}{}}{=} x_i^{t,fc_b}\otimes x_i^{t,fc_c}$. The authors then employ feature adaptation by means of Maximum Mean Discrepancy, between the source and target domains, on these fusion features.<br />
<br />
[[File:RTN-MMD-Block.png|right|thumb|The Maximum Mean Discrepancy Block (MMD) included in the RTN. The outputs of $fc_b$ and $fc_c$ are fused through a tensor product, and then passed through the MMD penalty, ensuring distributional similarity.]]<br />
<br />
==== Maximum Mean Discrepancy ==== <br />
The Maximum Mean Discrepancy (MMD) is a Kernel method that involves mapping to a Reproducing Kernel Hilbert Space (RKHS) [https://en.wikipedia.org/wiki/Reproducing_kernel_Hilbert_space]. Denote the RKHS $\mathcal{H}_K$ with a characteristic kernel $K$. We then define the '''mean embedding''' of a distribution $p$ in $\mathcal{H}_K$ to be the unique element $\mu_K(p)$ such that $\mathbf{E}_{x\sim p}f(x) = \langle f(x), \mu_K(p)\rangle_{\mathcal{H}_K}$ for all $f \in \mathcal{H}_K$. Now, if we take $\phi: \mathcal{X} \to \mathcal{H}_K$, then we can define the MMD between two distributions $p$ and $q$ as follows:<br />
<br />
<center><br />
<math display="block"><br />
d_k(p, q) \overset{\underset{\mathrm{def}}{}}{=} ||\mathbf{E}_{x\sim p}(\phi(x^s)) - \mathbf{E}_{x\sim q}(\phi(x^t))||_{\mathcal{H}_K}<br />
</math><br />
</center><br />
<br />
Effectively, the MMD will compute the self-similarity of $p$ and $q$, and subtract twice the cross-similarity between the distributions: $\widehat{\text{MMD}}^2 = \text{mean}(K_{pp}) + \text{mean}(K_{qq}) - 2\times\text{mean}(K_{pq})$. From here we can infer that $p$ and $q$ are equivalent distributions if and only if the $\text{MMD} = 0$. If we then wish to force two distributions to be similar, this becomes a minimization problem over the MMD.<br />
<br />
Two important notes:<br />
# The RKHS, and as such MMD, depend on the choice of the kernel;<br />
# Computing the MMD efficiently requires an unbiased estimate of the MMD (as outlined [https://arxiv.org/pdf/1502.02791.pdf]).<br />
<br />
==== MMD for Feature Adaptation in the RTN ====<br />
The authors wish to minimize the MMD between the fusion features outlined above derived from the source and target domains. Concretely this amounts to forcing the distribution of the abstract representation of the source domain $\mathcal{D}_s$ to be similar to the distribution of the abstract representation of the target domain $\mathcal{D}_t$. Performing this optimization over the fused features between the $fb_b$ and $fb_c$ forces each of those layers towards similar distributions.<br />
<br />
Practically this involves an additional penalty function given by the following:<br />
<br />
<center><br />
<math display="block"><br />
D_{\mathcal{L}}(\mathcal{D}_s, \mathcal{D}_t) = \sum_{i,j=1}^{n_s} \frac{k(z_i^s, z_j^s)}{n_s^2} + \sum_{i,j=1}^{n_t} \frac{k(z_i^t, z_j^t)}{n_t^2} + \sum_{i=1}^{n_s}\sum_{j=1}^{n_t} \frac{k(z_i^s, z_j^t)}{n_sn_t} <br />
</math><br />
</center><br />
<br />
Where the characteristic kernel $k(z, z')$ is the Gaussian kernel, defined on the vectorization of tensors, with bandwidth parameter $b$. That is: $k(z, z') = \exp(-||vec(z) - vec(z')||^2/b)$.<br />
<br />
=== Classifier Adaptation ===<br />
In traditional unsupervised domain adaptation there is a '''shared-classifier assumption''' which is made. In essence, if $f_s(x)$ represents the classifier on the source domain, and $f_t(x)$ represents the classifier on the target domain then this assumption simply states that $f_s = f_t$. While this may seem to be a reasonable assumption at first glance, it is problematic largely in that this is an assumption that is incredibly difficult to check. If it could be readily confirmed that the source and target classifiers could be shared, then the problem of domain adaptation would be largely trivialized. Instead, the authors here relax this assumption slightly. They postulate that instead of being equivalent, the source and target classifier differ by some perturbation function $\Delta f$. The general idea is that, by assuming $f_S(x) = f_T(x) + \Delta f(x)$, where $f_S$ and $f_T$ correspond to the source and target classifiers, pre-activation, and $\Delta f(x)$ is some residual function.<br />
<br />
The authors then suggest using residual blocks, as popularized by the ResNet framework [https://arxiv.org/pdf/1512.03385.pdf], to learn this residual function.<br />
<br />
[[File:Residual-Block-vs-DNN.png|thumb|left|A comparison of a standard Deep Neural Network block which is designed to fit a function H(x) compared to a residual block which fits H(x) as the sum of the input, x, and a learned residual function, F(X).]]<br />
==== Residual Networks Framework ==== <br />
A (Deep) Residual Network, as proposed initially in ResNet, employs residual blocks to assist in the learning process, and were a key component of being able to train extraordinarily deep networks. The Residual Network is comprised largely in the same manner as standard neural networks, with one key difference, namely the inclusion of residual blocks - sets of layers which aim to estimate a residual function in place of estimating the function itself. <br />
<br />
That is, if we wish to use a DNN to estimate some function $h(x)$, a residual block will decompose this to $h(x) = F(x) + x$. The layers are then used to learn $F(x)$, and after the layers which aim to learn this residual function, the input $x$ is recombined through element-wise addition, to form $h(x) = F(x) + x$. This was initially proposed as a manner to allow for deeper networks to be effectively trained, but has since used in novel contexts.<br />
<br />
==== Residual Blocks in the RTN ====<br />
[[File:RTN-Residual-Block.png|thumb|right|The Structure of the Residual Block in the RTN framework. The block relies on two additional dense layers following the target classifier in an attempt to learn the residual difference between the source and target classifiers.]] The authors leverage residual blocks for the purpose of classifier adaptation. Operating under the assumption that the source and target classifiers differ by an arbitrary perturbation function, $f(x)$, the authors add an additional set of densely connected layers which the source data will flow through. In particular, the authors take the $fc_c$ layer above as the desired target classifier. For the source data an additional set of layers ($fc-1$ and $fc-2$) are added following $fc_c$, which are connected as a residual block. The output of the classifier layer is then added back to the output of the residual block in order to form the source classifier.<br />
<br />
It is necessary to note that in this case the output from $fc_c$ passes the non-activated (i.e. pre-softmax activation) to the element-wise addition, the result of which is passed through the activation layer, yielding the source prediction. In the provided diagram, we have that $f_S(x)$ represents the non-activated output from the additive layer in the residual block; $f_T(x)$ represents the non-activated output from the target classifier; and $fc-1$/$fc-2$ are used to learn the perturbation function $\Delta f(x)$.<br />
<br />
==== Entropy Minimization ====<br />
In addition to the residual blocks, the authors make use of the '''entropy minimization principle''' [http://www.iro.umontreal.ca/~lisa/pointeurs/semi-supervised-entropy-nips2004.pdf] to further refine the classifier adaptation. In particular, by minimizing the entropy of the target classifier (or more correctly, the entropy of the class conditional distribution $f_j^t(x_i^t) = p(y_i^t = j \mid x_i^t; f_t)$), low-density separation between the classes is encouraged. '''Low-Density Separation''' is a concept used predominantly in semi-supervised learning, which in essence tries to draw class decision boundaries in regions where there are few data points (labeled or unlabeled). The above paper leverages an entropy regularization scheme to achieve the goal low-density separation goal; this is adopted here to the case of unsupervised domain adaptation.<br />
<br />
In practice this amounts to adding a further penalty based on the entropy of the class conditional distribution. In particular, if $H(\cdot)$ is defined to be the entropy function, such that $H(f_t(x_i^t)) = - \sum_{j=1}^c f_j^t(x_i^t)\log f_j^t(x_i^t)$, where $c$ is the number of classes and $f_j^t(x_i^t)$ represents the probability of predicting class $j$ for point $x_i^t$, then over the target domain $\mathcal{D}_t$ we define the entropy penalty to be:<br />
<br />
<center><br />
<math display="block"><br />
\frac{1}{n_t} \sum_{i=1}^{n_t} H(f_t(x_i^t))<br />
</math><br />
</center><br />
<br />
The combination of the residual learning and the entropy penalty, the authors hypothesize will enable effective classifier adaptation.<br />
<br />
=== Residual Transfer Network ===<br />
The combination of the MMD loss introduced in feature adaptation, the residual block introduced in classifier adaptation, and the application of the entropy minimization principle cumulates in the Residual Transfer Network proposed by the authors. The model will be optimized according to the following loss function, which combines the standard cross-entropy, MMD penalty, and entropy penalty:<br />
<br />
<center><br />
<math display="block"><br />
\min_{f_s = f_t + \Delta f} \underbrace{\left(\frac{1}{n_s} \sum_{i=1}^{n_s} L(f_s(x_i^s), y_i^s)\right)}_{\text{Typical Cross-Entropy}} + \underbrace{\frac{\gamma}{n_t}\left(\sum_{i=1}^{n_t} H(f_t(x_i^t)) \right)}_{\text{Target Entropy Minimization}} + \underbrace{\lambda\left(D_{\mathcal{L}}(\mathcal{D}_s, \mathcal{D}_t)\right)}_{\text{MMD Penalty}}<br />
</math><br />
</center><br />
<br />
Where we take $\gamma$ and $\lambda$ to be tradeoff parameters between the entropy penalty and the MMD penalty. As classifier adaptation proposed in this paper and feature adaptation studied in [5, 6] are tailored to adapt different layers of deep networks, they are expected to complement each other and to establish better performance.<br />
<br />
The full network, which is trained subject to the above optimization problem, thus takes on the following structure.<br />
<br />
[[File:rtn-full-paper-structure.png||center|alt=The Structure of the RTN]]<br />
<br />
== Experiments == <br />
<br />
=== Set-up ===<br />
The performance of RTN was jointly compared across two key data sets in the area of Unsupervised Domain Adaptation. Specifically, Office-31 (discussed in the introduction) and Office-Caltech (maintained by the same project group). Office-31 is comprised of images from 3 sources, Amazon ('''A'''), Webcam ('''W'''), and DSLR ('''D'''), of 31 different objects. Office-Caltech is derived by considering 10 classes common to both the Office-31 and the Caltech data sets, thus providing further adaptation possibilities. This provides 6 Transfer Tasks on the 31 classes of Office-31 ($\{(A,W), (A,D), (W,A), (W,D), (D,A), (D,W)\}$) and 12 Transfer Tasks on the 10 classes of Office-Caltech ($\{(A,W), (A,D), (A,C), (W,A), (W,D), (W,C), (D,A), (D,W), (D,C), (C,A), (C,W), (C,D)\}$).<br />
<br />
The authors then compare the results on the 18 different adaptation tasks against 6 other models. In order to determine the efficacy of the various contributions outlined in the paper they perform an ablation study, evaluating variants of the RTN. Specifically, they consider the RTN with only the MMD module ('''RTN (mmd)'''), the RTN with the MMD module and the entropy minimization ('''RTN (mmd+ent)'''), and the complete RTN ('''RTN (mmd+ent+res)'''). The experiments leverage all the labeled training data and compute accuracy across all unlabeled domain data. The parameters of the model (i.e. $\gamma$, and $\lambda$) are fixed based on a single validation point on the transfer task $\mathbf{A}\to\mathbf{W}$. These parameters are then maintained across all transfer tasks. <br />
<br />
As for specification details, the authors use mini-batch SGD, with momentum $0.9$, and with the learning rate adjusted based on $\eta_p = \frac{\eta_0}{(1 + \alpha p)^\beta}$, where $p$ indicates the portion of training completed (linear from $0$ to $1$), $\eta_0 = 0.01$, $\alpha = 10$ and $\beta = 0.75$, which was optimized for low error on the source. The MMD and entropy parameters, set as above, were maintained at $\lambda = 0.3$ and $\gamma - 0.3$.<br />
<br />
=== Results ===<br />
[[File:table-1-results.PNG|thumb|right|Results from the Office-31 Experiment]][[File:table-2-results.PNG|thumb|right|Results from the Office-Caltech Experiment]]<br />
In aggregate, the network outperformed all comparison methods, across all transfer tasks. Broadly speaking the network saw the largest increases in accuracy on the hard transfer tasks (for instance $\mathbf{A} \to \mathbf{C}$), where the source-domain discrepancy is large. The authors take this to mean that the proposed model learns "more adaptive classifiers and transferable features for safer domain adaptation." They further indicate that standard deep learning techniques (i.e. just AlexNet) perform similarly to standard shallow techniques (TCA and GFK). Deep-transfer methods which focus on feature adaptation perform significantly better than the standard methods. The proposed RTN, which adds in additional considerations for classifier adaptation, performs even better.<br />
<br />
In addition, the ablation study found a number of interesting results:<br />
# The RTN (mmd) outperforms DAN, which is founded on a similar method, but contains multiple MMD penalties (one for each layer instead of on a bottleneck), and is as such less computationally efficient;<br />
# The addition of the entropy penalty [RTN (mmd+ent)] provides significant marginal benefit over the previous RTN (mmd);<br />
# The full RTN [RTN (mmd+ent+res)] performs the best of all variants, by diminishing returns are seen over the addition of the entropy penalty.<br />
<br />
Overall the authors claim that the RTN (mmd+ent+res) is now regarded as state-of-the-art for unsupervised domain adaptation.<br />
<br />
=== Discussion ===<br />
[[File:t-sne-embeddings.png|thumb|left|t-SNE Embeddings Comparing the Performance of DAN and RTN]] <br />
[[File:mean-sd-layer-outputs.png|thumb|right|The Mean and Standard Deviations of the outputs from the Source Classifier, Target Classifier, and Residual Functions. As expected, the residual function provides a small, but non-zero, contribution.]] <br />
[[File:gamma-tradeoff.png|thumb|left|The accuracy of tests by varying the parameter $\gamma$. We first see an increase in accuracy up to an ideal point, before having the accuracy fall again.]]<br />
[[File:classifier-shift.png|thumb|right|The corresponding weights of the classifier layers, if trained on the labeled source and target data, exhibiting the differences which exist between the two classifiers in an ideal state. ]]<br />
<br />
==== Visualizing Predictions (Versus DAN) ====<br />
DAN uses a similar method for feature adaptation but neglects any attempt at classifier adaptation (i.e. it makes the shared-classifier assumption). In order to demonstrate that this leads to the worse performance, the authors provide images showing the t-SNE embeddings by DAN and RTN on the transfer task $\mathbf{A} \to \mathbf{W}$. The images show that the target categories are not well discriminated by the source classifier, suggesting a violation of the shared-classifier assumption. Conversely, the target classifier for the RTN exhibits better discrimination.<br />
<br />
==== Layer Responses and Classifier Shift ==== <br />
The authors further consider the mean and standard deviation of the outputs of $f_S(x)$, $f_T(x)$ and $\Delta f(x)$ to consider the relative contributions of the different components. As expected, $\Delta f(x)$ provides a small (though non-zero) contribution to the learned source classifier. This provides some merit to the idea of residual learning on the classifiers. <br />
<br />
In addition, the authors train classifiers on the source and target data, with labels present, and compare the realized weights. This is used to test how different the ideal weights are on separate classifiers. The results suggest that there is, in fact, a discrepancy between the classifiers, further motivating the use of tactics to avoid the shared-classifier assumption. <br />
<br />
==== Parameter Sensitivity ==== <br />
Lastly, the authors test the sensitivity of these results against the parameter $\gamma$. They run this test on $\mathbf{A}\to\mathbf{W}$ in addition to $\mathbf{C}\to\mathbf{W}$, varying the parameter from $0.01$ to $1.0$. They find that, on both tasks, the increase of the parameter initially improves accuracy, before seeing a drop-off.<br />
<br />
== Conclusion ==<br />
This paper presented a novel approach to unsupervised domain adaptation which relaxed assumptions made by previous models with regard to the shared nature of classifiers. Like previous models this proposed network leverages feature adaptation by matching the distributions of features across the domains. In addition, using a residual network and entropy minimization tactic, the target classifier is allowed to differ from the source classifier by implementing a new residual transfer module as the bridge. In particular, this approach allows for easy integration into existing networks, and can be implemented with any standard deep learning software.<br />
<br />
For follow-up considerations, the authors propose looking for adaptations which may be useful in the semi-supervised domain adaptation problem.<br />
<br />
== Critique ==<br />
While the paper presents a clear approach, which empirically attains great results on the desired tasks, I question the benefit to the residual block that is employed. The results of the ablation study seem to suggest that the majority of the benefits can be derived from using the MMD and Entropy penalties. The residual block appears to add marginal, perhaps insignificant contributions to the outcome. Despite this, the use of MMD loss is not novel, and the entropy loss is less well documented, and less thoroughly explored. Perhaps a different set of ablations would have indicated that the three parts, indeed, are equally effective (and the diminishing returns stems from stacking the three methods), but as it is presented, I question the utility of the final structure versus a less complicated, less novel approach.<br />
<br />
==References==<br />
# https://en.wikipedia.org/wiki/Domain_adaptation<br />
# https://people.eecs.berkeley.edu/~jhoffman/domainadapt/<br />
# Glorot, Xavier, Antoine Bordes, and Yoshua Bengio. "Domain adaptation for large-scale sentiment classification: A deep learning approach." Proceedings of the 28th international conference on machine learning (ICML-11). 2011.<br />
# Tzeng, Eric, et al. "Deep domain confusion: Maximizing for domain invariance." arXiv preprint arXiv:1412.3474 (2014).<br />
# Long, Mingsheng, et al. "Learning transferable features with deep adaptation networks." International Conference on Machine Learning. 2015.<br />
# Ganin, Yaroslav, and Victor Lempitsky. "Unsupervised domain adaptation by backpropagation." International Conference on Machine Learning. 2015.<br />
# Tzeng, Eric, et al. "Simultaneous deep transfer across domains and tasks." Proceedings of the IEEE International Conference on Computer Vision. 2015.<br />
# He, Kaiming, et al. "Deep residual learning for image recognition." Proceedings of the IEEE conference on computer vision and pattern recognition. 2016.<br />
# Yang, Jun, Rong Yan, and Alexander G. Hauptmann. "Cross-domain video concept detection using adaptive svms." Proceedings of the 15th ACM international conference on Multimedia. ACM, 2007.<br />
# Duan, Lixin, et al. "Domain adaptation from multiple sources via auxiliary classifiers." Proceedings of the 26th Annual International Conference on Machine Learning. ACM, 2009.<br />
# Duan, Lixin, et al. "Visual event recognition in videos by learning from web data." IEEE Transactions on Pattern Analysis and Machine Intelligence 34.9 (2012): 1667-1680.<br />
# http://vision.stanford.edu/teaching/cs231b_spring1415/slides/alexnet_tugce_kyunghee.pdf<br />
# https://en.wikipedia.org/wiki/Reproducing_kernel_Hilbert_space<br />
#Long, Mingsheng, et al. "Learning transferable features with deep adaptation networks." International Conference on Machine Learning. 2015.<br />
#He, Kaiming, et al. "Deep residual learning for image recognition." Proceedings of the IEEE conference on computer vision and pattern recognition. 2016.<br />
# Grandvalet, Yves, and Yoshua Bengio. "Semi-supervised learning by entropy minimization." Advances in neural information processing systems. 2005.<br />
<br />
Expert review from the NIPS community can be found in https://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/99.html.<br />
<br />
Implementation Example: https://github.com/thuml/Xlearn</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Model-Agnostic_Meta-Learning_for_Fast_Adaptation_of_Deep_Networks&diff=31011Model-Agnostic Meta-Learning for Fast Adaptation of Deep Networks2017-11-21T17:27:42Z<p>Nkadiyar: </p>
<hr />
<div>='''Introduction & Background'''=<br />
Learning quickly is a hallmark of human intelligence, whether it involves recognizing objects from a few examples or quickly learning new skills after just minutes of experience. In this work, we propose a meta-learning algorithm that is general and model-agnostic, in the sense that it can be directly applied to any learning problem and model that is trained with a gradient descent procedure. Our focus is on deep neural network models, but we illustrate how our approach can easily handle different architectures and different problem settings, including classification, regression, and policy gradient reinforcement learning, with minimal modification. Unlike prior meta-learning methods that learn an update function or learning rule (Schmidhuber, 1987; Bengio et al., 1992; Andrychowicz et al., 2016; Ravi & Larochelle, 2017), this algorithm does not expand the number of learned parameters nor place constraints on the model architecture (e.g. by requiring a recurrent model (Santoro et al., 2016) or a Siamese network (Koch, 2015)), and it can be readily combined with fully connected, convolutional, or recurrent neural networks. It can also be used with a variety of loss functions, including differentiable supervised losses and nondifferentiable reinforcement learning objectives.<br />
<br />
The primary contribution of this work is a simple model and task-agnostic algorithm for meta-learning that trains a model’s parameters such that a small number of gradient updates will lead to fast learning on a new task. The paper shows the effectiveness of the proposed algorithm in different domains, including classification, regression, and reinforcement learning problems.<br />
<br />
='''Model-Agnostic Meta Learning (MAML)'''=<br />
The goal of the proposed model is rapid adaptation. This setting is usually formalized as few-shot learning.<br />
<br />
=== Problem set-up ===<br />
The goal of few-shot meta-learning is to train a model that can quickly adapt to a new task using only a few data points and training iterations. To do so. the model is trained during a meta-learning phase on a set of tasks, such that it can then be adapted to a new task using only a small number of parameter updates. In effect, the meta-learning problem treats entire tasks as training examples. <br />
<br />
Let us consider a model denoted by $f$, that maps the observation $\mathbf{x}$ into the output variable $a$. During meta-learning, the model is trained to be able to adapt to a large or infinite number of tasks. <br />
<br />
Let us consider a generic notion of task as below. Each task $\mathcal{T} = \{\mathcal{L}(\mathbf{x}_1.a_1,\mathbf{x}_2,a_2,..., \mathbf{x}_H,a_H), q(\mathbf{x}_1),q(\mathbf{x}_{t+1}|\mathbf{x}_t,a_t),H \}$, consists of a loss function $\mathcal{L}$, a distribution over initial observations $q(\mathbf{x}_1)$, a transition distribution $q(\mathbf{x}_{t+1}|\mathbf{x}_t)$, and an episode length $H$. In i.i.d. supervised learning problems,<br />
the length $H =1$. The model may generate samples of length $H$ by choosing an output at at each time $t$. The cost $\mathcal{L}$ provides a task-specific feedback, which is defined based on the nature of the problem. <br />
<br />
A distribution over tasks is denoted by $p(\mathcal{T})$. In the K-shot learning setting, the model is trained to learn a new task $\mathcal{T}_i$ drawn from $p(\mathcal{T})$ from only K samples drawn from $q_i$ and feedback $\mathcal{L}_{\mathcal{T}_i}$ generated by $\mathcal{T}_i$. During meta-training, a task $\mathcal{T}_i$ is sampled from $p(\mathcal{T})$, the model is trained with K samples and feedback from the corresponding loss LTi from Ti, and then tested on new samples from Ti. The model f is then improved by considering how the test error on new data from $q_i$ changes with respect to the parameters. In effect, the test error on sampled tasks $\mathcal{T}_i$ serves as the training error of the meta-learning process. At the end of meta-training, new tasks are sampled from $p(\mathcal{T})$, and meta-performance is measured by the model’s performance after learning from K samples.<br />
<br />
=== MAML Algorithm ===<br />
[[File:model.png|200px|right|thumb|Figure 1: Diagram of the MAML algorithm]]<br />
The paper proposes a method that can learn the parameters of any standard model via meta-learning in such a way as to prepare that model for fast adaptation. The intuition behind this approach is that some internal representations are more transferrable than others. Since the model will be fine-tuned using a gradient-based learning rule on a new task, we will aim to learn a model in such a way that this gradient-based learning rule can make rapid progress on new tasks drawn from $p(\mathcal{T})$, without overfitting. In effect, we will aim to find model parameters that are sensitive to changes in the task, such that small changes in the parameters will produce large improvements on the loss function of any task drawn from $p(\mathcal{T})$, see Fig 1.<br />
<br />
Note that there is no assumption about the form of the model. Only assumption is that it is parameterized by a vector of parameters $\theta$, and the loss is smooth so that the parameters can be leaned using gradient-based techniques. Formally lets assume that the model is denoted by $f_{\theta}$. When adapting<br />
to a new task $\mathcal{T}_i $, the model’s parameters $\theta$ become $\theta_i'$. In our method, the updated parameter vector $\theta_i'$ is computed using one or more gradient descent updates on task $\mathcal{T}_i $. For example, when using one gradient update:<br />
<br />
$$<br />
\theta_i ' = \theta - \alpha \nabla_{\theta \mathcal{L}_{\mathcal{T}_i}}(f_{\theta}).<br />
$$<br />
<br />
Here $\alpha$ is a the learning rate of each task and considered as a hyperparameter. They consider a single step of update for the rest of the paper, for the sake of the simplicity. <br />
<br />
The model parameters are trained by optimizing for the performance<br />
of $f_{\theta_i'}$ with respect to $\theta$ across tasks sampled from $p(\mathcal{T})$. More concretely, the meta-objective is as follows: <br />
<br />
$$<br />
\min_{\theta} \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta_i'}) = \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta - \alpha \nabla_{\theta \mathcal{L}_{\mathcal{T}_i}}(f_{\theta})})<br />
$$<br />
<br />
Note that the meta-optimization is performed over the model parameters $\theta$, whereas the objective is computed using the updated model parameters $\theta'$. The model aims to optimize the model parameters such that one or a small number of gradient step on a new task will produce maximally effective behavior on that task. <br />
<br />
Therefore the meta-learning across the tasks is performed via stochastic gradient descent (SGD), such that the model parameters $\theta$ are updated as:<br />
<br />
$$<br />
\theta \gets \theta - \beta \nabla_{\theta } \sum \limits_{\mathcal{T}_i \sim p(\mathcal{T})} \mathcal{L}_{\mathcal{T}_i} (f_{\theta_i'})<br />
$$<br />
where $\beta$ is the meta step size. Outline of the algorithm is shown in Algorithm 1. <br />
[[File:ershad_alg1.png|500px|center|thumb]]<br />
<br />
The MAML meta-gradient update involves a gradient through a gradient. Computationally, this requires an additional backward pass through f to compute Hessian-vector products, which is supported by standard deep learning libraries such as TensorFlow.<br />
<br />
='''Different Types of MAML'''=<br />
In this section the MAML algorithm is discussed for different supervised learning and reinforcement learning tasks. The differences between each of these tasks are in their loss function and the way the data is generated. In general, this method does not require additional model parameters nor using any additional meta-learner to learn the update of parameters. Compared to other approaches that tend to “learn to compare new examples in a learned metric space using e.g. Siamese networks or recurrence with attention mechanisms”, the proposed method can be generalized to any other problems including classification, regression and reinforcement learning. <br />
<br />
=== Supervised Regression and Classification ===<br />
Few-shot learning is well-studied in this field. For these two types of tasks the horizon $H$ is equal to 1, since the data points are generated i.i.d. <br />
<br />
Although any common classification and regression objectives can be used as the loss, the paper uses the following losses for these two tasks. <br />
<br />
Regression : For regression we use the mean-square error (MSE):<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i} (f_{\theta}) = \sum \limits_{\mathbf{x}^{(j)}, \mathcal{y}^{(j)} \sim \mathcal{T}_i} \parallel f_{\theta} (\mathbf{x}^{(j)} - \mathbf{y}^{(j)})\parallel_2^2<br />
$$<br />
<br />
where $\mathbf{x}^{(j)}$ and $\mathbf{y}^{()j}$ are the input/output pair sampled from task $\mathcal{T}_i$. In K-shot regression tasks, K input/output pairs are provided for learning for each task. <br />
<br />
Classification: For classification we use the cross entropy loss:<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i} (f_{\theta}) = \sum \limits_{\mathbf{x}^{(j)}, \mathcal{y}^{(j)} \sim \mathcal{T}_i} \mathbf{y}^{(j)} \log f_{\theta}(\mathbf{x}^{(j)}) + (1-\mathbf{y}^{(j)}) \log (1-f_{\theta}(\mathbf{x}^{(j)}))<br />
$$<br />
<br />
According to the conventional terminology, K-shot classification tasks use K input/output pairs from each class, for a total of $NK$ data points for N-way classification.<br />
<br />
Given a distribution over tasks, these loss functions can be directly inserted into the equations in the previous section to perform meta-learning, as detailed in Algorithm 2.<br />
[[File:ershad_alg2.png|500px|center|thumb]]<br />
<br />
=== Reinforcement Learning ===<br />
In reinforcement learning (RL), the goal of few-shot meta learning is to enable an agent to quickly acquire a policy for a new test task using only a small amount of experience in the test setting. A new task might involve achieving a new goal or succeeding on a previously trained goal in a new environment. For example an agent may learn how to navigate mazes very quickly so that, when faced with a new maze, it can determine how to reliably reach the exit with only a few samples.<br />
<br />
Each RL task $\mathcal{T}_i$ contains an initial state distribution $q_i(\mathbf{x}_1)$ and a transition distribution $q_i(\mathbf{x}_{t+1}|\mathbf{x}_t,a_t)$ , and the loss $\mathcal{L}_{\mathcal{T}_i}$ corresponds to the (negative) reward function $R$. The entire task is therefore a Markov decision process (MDP) with horizon H, where the learner is allowed to query a limited number of sample trajectories for few-shot learning. Any aspect of the MDP may change across tasks in $p(\mathcal{T})$. The model being learned, $f_{\theta}$, is a policy that maps from states $\mathbf{x}_t$ to a distribution over actions $a_t$ at each timestep $t \in \{1,...,H\}$. The loss for task $\mathcal{T}_i$ and model $f_{\theta}$ takes the form<br />
<br />
$$<br />
\mathcal{L}_{\mathcal{T}_i}(f_{\theta}) = -\mathbb{E}_{\mathbf{x}_t,a_t \sim f_{\theta},q_{\mathcal{T}_i}} \big [\sum_{t=1}^H R_i(\mathbf{x}_t,a_t)\big ]<br />
$$<br />
<br />
<br />
In K-shot reinforcement learning, K rollouts from $f_{\theta}$ and task $\mathcal{T}_i$, $(\mathbf{x}_1,a_1,...,\mathbf{x}_H)$, and the corresponding rewards $ R(\mathbf{x}_t,a_t)$, may be used for adaptation on a new task $\mathcal{T}_i$.<br />
<br />
Since the expected reward is generally not differentiable due to unknown dynamics, we use policy gradient methods to estimate the gradient both for the model gradient update(s) and the meta-optimization. Since policy gradients are an on-policy algorithm, each additional gradient step during the adaptation of $f_{\theta}$ requires new samples from the current policy $f_{\theta_i}$ . We detail the algorithm in Algorithm 3, which has the same structure as Algorithm 2 but also which requires sampling trajectories from the environment corresponding to task $\mathcal{T}_i$ in steps 5 and 8.<br />
[[File:ershad_alg3.png|500px|center|thumb]]<br />
<br />
='''Experiments'''=<br />
<br />
=== Regression ===<br />
We start with a simple regression problem that illustrates the basic principles of MAML. Each task involves regressing from the input to the output of a sine wave, where the amplitude and phase of the sinusoid are varied between tasks. Thus, $p(\mathcal{T})$ is continuous, and the input and output both have a dimensionality of 1. During training and testing, datapoints are sampled uniformly. The loss is the mean-squared error between the prediction and true value. The regressor is a neural network model with 2 hidden layers of size 40 with ReLU nonlinearities. When training with MAML, we use one gradient update with K = 10 examples with a fixed step size 0.01, and use Adam as the metaoptimizer [2]. The baselines are likewise trained with Adam. To evaluate performance, we fine-tune a single meta-learned model on varying numbers of K examples, and compare performance to two baselines: (a) pre-training on all of the tasks, which entails training a network to regress to random sinusoid functions and then, at test-time, fine-tuning with gradient descent on the K provided points, using an automatically tuned step size, and (b) an oracle which receives the true amplitude and phase as input.<br />
<br />
We evaluate performance by fine-tuning the model learned by MAML and the pre-trained model on $K = \{ 5,10,20 \}$ datapoints. During fine-tuning, each gradient step is computed using the same $K$ datapoints. Results are shown in Fig 2.<br />
<br />
<br />
[[File:ershad_results1.png|500px|center|thumb|Figure 2: Few-shot adaptation for the simple regression task. Left: Note that MAML is able to estimate parts of the curve where there are no datapoints, indicating that the model has learned about the periodic structure of sine waves. Right: Fine-tuning of a model pre-trained on the same distribution of tasks without MAML, with a tuned step size. Due to the often contradictory outputs on the pre-training tasks, this model is unable to recover a suitable representation and fails to extrapolate from the small number of test-time samples.]]<br />
<br />
=== Classification ===<br />
<br />
For classification evaluation, Omniglot and MiniImagenet datasets are used. The Omniglot dataset consists of 20 instances of 1623 characters from 50 different alphabets. <br />
<br />
The experiment involves fast learning of N-way classification with 1 or 5 shots. The problem of N-way classification is set up as follows: select N unseen classes, provide the model with K different instances of each of the N classes, and evaluate the model’s ability to classify new instances within the N classes. For Omniglot, 1200 characters are selected randomly for training, irrespective of alphabet, and use the remaining for testing. The Omniglot dataset is augmented with rotations by multiples of 90 degrees.<br />
<br />
The model follows the same architecture as the embedding function that has 4 modules with a 3-by-3 convolutions and 64 filters, followed by batch normalization, a ReLU nonlinearity, and 2-by-2 max-pooling. The Omniglot images are downsampled to 28-by-28, so the dimensionality of the last hidden layer is 64. The last layer is fed into a softmax. For Omniglot, strided convolutions is used instead of max-pooling. For MiniImagenet, 32 filters per layer are used to reduce overfitting. In order to also provide a fair comparison against memory-augmented neural networks [3] and to test the flexibility of MAML, the results for a non-convolutional network are also provided. <br />
<br />
[[File:ershad_results2.png|500px|center|thumb|Table 1: Few-shot classification on held-out Omniglot characters (top) and the MiniImagenet test set (bottom). MAML achieves results that are comparable to or outperform state-of-the-art convolutional and recurrent models. Siamese nets, matching nets, and the memory module approaches are all specific to classification, and are not directly applicable to regression or RL scenarios. The $\pm$ shows 95% confidence intervals over tasks. ]]<br />
<br />
=== Reinforcement Learning ===<br />
Several simulated continuous control environments are used for RL evaluation. In all of the domain, the MAML model is a neural network policy with two hidden layers of size 100, and ReLU activations. The gradient updates are computed using vanilla policy gradient and trust-region policy optimization (TRPO) is used as the meta-optimizer.<br />
<br />
In order to avoid computing third derivatives, finite differences are computed to <br />
compute the Hessian-vector products for TRPO. For both learning and meta-learning updates, we use the standard linear feature baseline proposed by [4], which is fitted separately at each iteration for each sampled task in the batch. <br />
<br />
Three baseline models for the comparison are: <br />
(a) pretraining one policy on all of the tasks and then fine-tuning<br />
(b) training a policy from randomly initialized weights<br />
(c) an oracle policy which receives the parameters of the task as input, which for the tasks below corresponds to a goal position, goal direction, or goal velocity for the agent. <br />
<br />
The baseline models of (a) and (b) are fine-tuned with gradient descent with a manually tuned step size.<br />
<br />
2D Navigation: In the first meta-RL experiment, the authors study a set of tasks where a point agent must move to different goal positions in 2D, randomly chosen for each task within a unit square. The observation is the current 2D position, and actions correspond to velocity commands clipped to be in the range [-0.1; 0.1]. The reward is the negative squared distance to the goal, and episodes terminate when the agent is within 0:01 of the goal or at the horizon ofH = 100. The policy was trained with MAML <br />
to maximize performance after 1 policy gradient update using 20 trajectories. They compare adaptation to a new task with up to 4 gradient updates, each with 40 samples. Results are shown in Fig. 3.<br />
<br />
[[File:ershad_results3.png|500px|center|thumb|Figure 3: Top: quantitative results from 2D navigation task, Bottom: qualitative comparison between model learned with MAML and with fine-tuning from a pretrained network ]]<br />
<br />
Locomotion. To study how well MAML can scale to more complex deep RL problems, we also study adaptation on high-dimensional locomotion tasks with the MuJoCo simulator [5]. The tasks require two simulated robots – a planar cheetah and a 3D quadruped (the “ant”) – to run in a particular direction or at a particular velocity. In the goal velocity experiments, the reward is the negative absolute value between the current velocity of the agent and a goal, which is chosen uniformly at random between 0 and 2 for the cheetah and between 0 and 3 for the ant. In the goal direction experiments, the reward is the magnitude of the velocity in either the forward or backward direction, chosen at random for each task in p(T ). The horizon is H = 200, with 20 rollouts per gradient step for all problems except the ant forward/backward task, which used 40 rollouts per step. The results in Figure 5 show that MAML learns a model that can quickly adapt its velocity and direction with even <br />
just a single gradient update, and continues to improve with more gradient steps. The results also show that, on these challenging tasks, the MAML initialization substantially outperforms random initialization and pretraining.<br />
[[File:ershad_results4.png|500px|center|thumb|Figure 4: Reinforcement learning results for the half-cheetah and ant locomotion tasks, with the tasks shown on the far right. ]]<br />
<br />
A conceptual method to achieve fast adaptation in language modeling tasks ( not been experimented on by the authors) would be to explore methods of attaching an Attention Kernel which results in a simple and differentiable loss. It has been implemented in One-Shot Language Modeling along with state-of-the-art improvements in one-shot learning on Imagenet and Omniglot [7].<br />
<br />
='''Conclusion'''=<br />
<br />
The paper introduced a meta-learning method based on learning easily adaptable model parameters through gradient descent. The approach has a number of benefits. It is simple and does not introduce any learned parameters for meta-learning. It can be combined with any model representation that is amenable to gradient-based training, and any differentiable objective, including classification, regression, and reinforcement learning. Lastly, since the method merely produces a weight initialization, adaptation can be performed with any amount of data and any number of gradient steps, though it demonstrates state-of-the-art results on classification with only one or five examples per class. The authors also show that the method can adapt an RL agent using policy gradients and a very modest amount of experience.<br />
<br />
='''References'''=<br />
# Schmidhuber, J¨urgen. Learning to control fast-weight memories: An alternative to dynamic recurrent networks. Neural Computation, 1992.<br />
# Lake, Brenden M, Salakhutdinov, Ruslan, Gross, Jason, and Tenenbaum, Joshua B. One shot learning of simple visual concepts. In Conference of the Cognitive Science Society (CogSci), 2011.<br />
# Santoro, Adam, Bartunov, Sergey, Botvinick, Matthew, Wierstra, Daan, and Lillicrap, Timothy. Meta-learning with memory-augmented neural networks. In International Conference on Machine Learning (ICML), 2016.<br />
# Duan, Yan, Chen, Xi, Houthooft, Rein, Schulman, John, and Abbeel, Pieter. Benchmarking deep reinforcement learning for continuous control. In International Conference on Machine Learning (ICML), 2016.<br />
# Todorov, Emanuel, Erez, Tom, and Tassa, Yuval. Mujoco: A physics engine for model-based control. In International Conference on Intelligent Robots and Systems (IROS), 2012.<br />
# Videos the learned policies can be found in https://sites.google.com/view/maml.<br />
# Oriol Vinyals, Charles Blundell, Timothy Lillicrap, Koray Kavukcuoglu, Daan Wierstra. "Matching Networks for One Shot Learning". arXiv:1606.04080 [cs.LG]<br />
<br />
Implementation Example: https://github.com/cbfinn/maml</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Exploration_via_Bootstrapped_DQN&diff=31010Deep Exploration via Bootstrapped DQN2017-11-21T17:23:48Z<p>Nkadiyar: </p>
<hr />
<div>== Details ==<br />
<br />
'''Title''': Deep Exploration via Bootstrapped DQN<br />
<br />
'''Authors''': Ian Osband {1,2}, Charles Blundell {2}, Alexander Pritzel {2}, Benjamin Van Roy {1}<br />
<br />
'''Organisations''':<br />
# Stanford University<br />
# Google Deepmind<br />
<br />
'''Conference''': NIPS 2016<br />
<br />
'''URL''': [https://papers.nips.cc/paper/6501-deep-exploration-via-bootstrapped-dqn papers.nips.cc]<br />
<br />
'''Online code sources'''<br />
* [https://github.com/iassael/torch-bootstrapped-dqn github.com/iassael/torch-bootstrapped-dqn]<br />
<br />
This summary contains background knowledge from Section 2-7 (except Section 5). Feel free to skip if you already know.<br />
<br />
== Intro to Reinforcement Learning ==<br />
<br />
In reinforcement learning, an agent interacts with an environment with the goal to maximize its long term reward. A common application of reinforcement learning is to the [https://en.wikipedia.org/wiki/Multi-armed_bandit multi armed bandit problem]. In a multi armed bandit problem, there is a gambler and there are $n$ slot machines, and the gambler can choose to play any specific slot machine at any time. All the slot machines have their own probability distributions by which they churn out rewards, but this is unknown to the gambler. So the question is, how can the gambler learn how to get the maximum long term reward?<br />
<br />
There are two things the gambler can do at any instance: either he can try a new slot machine, or he can play the slot machine he has tried before (and he knows he will get some reward). However, even though trying a new slot machine feels like it would bring less reward to the gambler, it is possible that the gambler finds out a new slot machine that gives a better reward than the current best slot machine. This is the dilemma of '''exploration vs exploitation'''. Trying out a new slot machine is '''exploration''', while redoing the best move so far is '''exploiting''' the currently understood perception of the reward.<br />
<br />
[[File:multiarmedbandit.jpg|thumb|Source: [https://blogs.mathworks.com/images/loren/2016/multiarmedbandit.jpg blogs.mathworks.com]]]<br />
<br />
There are many strategies to approach this '''exploration-exploitation dilemma'''. Some [https://web.stanford.edu/class/msande338/lec9.pdf common strategies] for optimizing in an exploration-exploitation setting are Random Walk, Curiosity-Driven Exploration, and Thompson Sampling. A lot of these approaches are provably efficient, but assume that the state space is not very large. For instance, the approach called Curiosity-Driven Exploration aims to take actions that lead to immediate additional information. This requires the model to search “every possible cell in the grid” which is not desirable if state space is very large. Strategies for large state spaces often just either ignore exploration, or do something naive like $\epsilon$-greedy, where you exploit with $1-\epsilon$ probability and explore "randomly" in rest of the cases. The general idea to tackle large or continuous state spaces is by value function approximation. An empirically tested strategy is Value Function Approximation using Fourier Basis [16]. It has also proven to perform well compared to radial basis functions and the polynomial basis, which are the two most popular fixed bases for linear value function approximation. <br />
<br />
This paper tries to use a Thompson sampling like approach to make decisions.<br />
<br />
== Thompson Sampling<sup>[[#References|[1]]]</sup> ==<br />
<br />
In Thompson sampling, our goal is to reach a belief that resembles the truth. Let's consider a case of coin tosses (2-armed bandit). Suppose we want to be able to reach a satisfactory pdf for $\mathbb{P}_h$ (heads). Assuming that this is a Bernoulli bandit problem, i.e. the rewards are $0$ or $1$, we can start off with $\mathbb{P}_h^{(0)}=\beta(1,1)$. The $\beta(x,y)$ distribution is a very good choice for a possible pdf because it works well for Bernoulli rewards. Further $\beta(1,1)$ is the uniform distribution $\mathbb{N}(0,1)$.<br />
<br />
Now, at every iteration $t$, we observe the reward $R^{(t)}$ and try to make our belief close to the truth by doing a Bayesian computation. Assuming $p$ is the probability of getting a heads,<br />
<br />
$$<br />
\begin{align*}<br />
\mathbb{P}(R|D) &\propto \mathbb{P}(D|R) \cdot \mathbb{P}(R) \\<br />
\mathbb{P}_h^{(t+1)}&\propto \mbox{likelihood}\cdot\mbox{prior} \\<br />
&\propto p^{R^{(t)}}(1-p)^{R^{(t)}} \cdot \mathbb{P}_h^{(t)} \\<br />
&\propto p^{R^{(t)}}(1-p)^{R^{(t)}} \cdot \beta(x_t, y_t) \\<br />
&\propto p^{R^{(t)}}(1-p)^{R^{(t)}} \cdot p^{x_t-1}(1-p)^{y_t-1} \\<br />
&\propto p^{x_t+R^{(t)}-1}(1-p)^{y_t+R^{(t)}-1} \\<br />
&\propto \beta(x_t+R^{(t)}, y_t+R^{(t)})<br />
\end{align*}<br />
$$<br />
<br />
[[File:thompson sampling coin example.png|thumb||||600px|Source: [https://www.quora.com/What-is-Thompson-sampling-in-laymans-terms Quora]]]<br />
<br />
This means that with successive sampling, our belief can become better at approximating the truth. There are similar update rules if we use a non Bernoulli setting, say, Gaussian. In the Gaussian case, we start with $\mathbb{P}_h^{(0)}=\mathbb{N}(0,1)$ and given that $\mathbb{P}_h^{(t)}\propto\mathbb{N}(\mu, \sigma)$ it is possible to show that the update rule looks like<br />
<br />
$$<br />
\mathbb{P}_h^{(t+1)} \propto \mathbb{N}\bigg(\frac{t\mu+R^{(t)}}{t+1},\frac{\sigma}{\sigma+1}\bigg)<br />
$$<br />
<br />
=== How can we use this in reinforcement learning? ===<br />
<br />
We can use this idea to decide when to explore and when to exploit. We start with an initial belief, choose an action, observe the reward and based on the kind of reward, we update our belief about what action to choose next.<br />
<br />
== Bootstrapping <sup>[[#References|[2,3]]]</sup> ==<br />
<br />
This idea may be unfamiliar to some people, so I thought it would be a good idea to include this. In statistics, bootstrapping is a method to generate new samples from a given sample. Suppose that we have a given population, and we want to study a measure $\theta$. So, we just find $n$ sample points (sample $\{D_i\}_{i=1}^n$), calculate this measure $\hat{\theta}$ for these $n$ points, and make our inference. <br />
<br />
If we later wish to find a better bound on $\hat{\theta}$, i.e. suppose we want to say that $\delta_1 \leq \hat{\theta} \leq \delta_2$ with a confidence of $c$, then we can use bootstrapping for this.<br />
<br />
Using bootstrapping, we can create a new sample $\{D'_i\}_{i=1}^{n'}$ by '''randomly sampling $n'$ times from $D$, with replacement'''. So, if $D=\{1,2,3,4\}$, a $D'$ of size $n'=10$ could be $\{1,4,4,3,2,2,2,1,3,4\}$. We do this a sufficient $k$ number of times, calculate $\hat{\theta}$ each time, and thus get a distribution $\{\hat{\theta}_i\}_{i=1}^k$. Now, we can choose the $100\cdot c$<sup>th</sup> and $100\cdot(1-c)$<sup>th</sup> percentile of this distribution, (let them be $\hat{\theta}_\alpha$ and $\hat{\theta}_\beta$ respectively) and say<br />
<br />
$$\hat{\theta}_\alpha \leq \hat{\theta} \leq \hat{\theta}_\beta, \mbox{with confidence }c$$<br />
<br />
== Why choose bootstrap and not dropout? ==<br />
<br />
There is previous work<sup>[[#References|[4]]]</sup> that establishes dropout as a good way to train NNs on a posterior such that the trained NN works like a function approximator that is close to the actual posterior. But, there are several problems with the predictions of this trained NN. The figures below are from the appendix of this paper. The left image is the NN trained by the authors of this paper on a sample noisy distribution and the right image is from the accompanying web demo from [[#References|[4]]], where the authors of [[#References|[4]]] show that their NN converges around the mean with a good confidence.<br />
<br />
[[File:dropout_results.png|thumb||center||700px|Source: this paper's appendix]]<br />
<br />
According to the authors of this paper,<br />
# Even though [[#References|[4]]] says that dropout converges arond the mean, their experiment actually behaves weirdly around a reasonable point like $x=0.75$. They think that this happens because dropout only affects the region local to the original data.<br />
# Samples from the NN trained on the original data do not look like a reasonable posterior (very spiky).<br />
# The trained NN collapses to zero uncertainty at the data points from the original data.<br />
<br />
== Q Learning and Deep Q Networks <sup>[[#References|[5]]]</sup> ==<br />
<br />
At any point of time, our rewards dictate what our actions should be. Also, in general, we want good long term rewards. For example, if we are playing a first person shooter game, it is a good idea to go out of cover to kill an enemy, even if some health is lost. Similarly, in reinforcement learning, we want to maximize our long term reward. So if at each time $t$, the reward is $r_t$, then a naive way is to say we want to maximise<br />
<br />
$$<br />
R_t = \sum_{i=0}^{\infty}r_t<br />
$$<br />
<br />
But, this reward is unbounded. So technically it could tend to $\infty$ in a lot of the cases. This is why we use a '''discounted reward'''.<br />
<br />
$$<br />
R_t = \sum_{i=0}^{\infty}\gamma^t r_t<br />
$$<br />
<br />
Here, we take $0\leq \gamma \lt 1$. So, what this means is that we value our current reward the most ($r_0$ has a coefficient of $1$), but we also consider the future possible rewards. So if we had two choices: get $+4$ now and $0$ at all other timesteps, or get $-2$ now and $+2$ after $3$ timesteps for $20$ timesteps, we choose the latter ($\gamma=0.9$). This is because $(+4) < (-2)+0.9^3(2+0.9\cdot2+\cdots+0.9^{19}\cdot2)$.<br />
<br />
<br />
A '''policy''' $\pi: \mathbb{S} \rightarrow \mathbb{A}$ is just a function that tells us what action to take in a given state $s\in \mathbb{S}$. Our goal is to find the best policy $\pi^*$ that maximises the reward from a given state $s$. So, a '''value function''' is defined from $s$ (which the agent is in, at timestep $t$) and following the policy $\pi$ as $V^\pi(s) = \mathbb{E}[R_t]$. The optimal value function is then simply<br />
<br />
$$<br />
V^*(s)=\displaystyle\max_{\pi}V^\pi(s)<br />
$$<br />
<br />
For convenience however, it is better to work with the '''Q function''' $Q: \mathbb{S}\times\mathbb{A} \rightarrow \mathbb{R}$. $Q$ is defined similarly as $V$. It is the expected return after taking an action $a$ in the given state $s$. So, $Q^\pi(s,a)=\mathbb{E}[R_t|s,a]$. The optimal $Q$ function is<br />
<br />
$$<br />
Q^*(s,a)=\displaystyle\max_{\pi}Q^\pi(s,a)<br />
$$<br />
<br />
Suppose that we know $Q^*$. Then, if we know that we are supposed to start at $s$ and take an action $a$ right now, what is the best course of action from the next time step? We just choose the optimal action $a'$ at the next state $s'$ that we reach. The optimal action $a'$ at state $s'$ is simply the argument $a_x$ that maximises our $Q^*(s',\cdot)$.<br />
<br />
$$<br />
a'=\displaystyle\arg\max_{a_x} Q^*(s',a_x)<br />
$$<br />
<br />
So, our best expected reward from $s$ taking action $a$ is $\mathbb{E}[r_t+\gamma\mathbb{E}[R_{t+1}]]$. This is known as the '''Bellman equation''' in optimal control problem (By the way, its continuous form is called '''Hamilton-Jacobi-Bellman equation''' or HJB equation, which is a very important partial differential equation):<br />
<br />
$$<br />
Q^*(s,a)=\mathbb{E}[r_t+\gamma \displaystyle\max_{a_x} Q^*(s',a_x)]<br />
$$<br />
<br />
In Q learning, we use a deep neural network with weights $\theta$ as a function approximator for $Q^*$, since Bellman equation is indeed a non-linear PDE and very difficult to solve numerically. The '''naive way''' to do this is to design a deep neural network that takes as input the state $s$ and action $a$, and produces an approximation to $Q^*$. <br />
<br />
* Suppose our neural net weights are $\theta_i$ at iteration $i$.<br />
* We want to train our neural net on the case when we are at $s$, take action $a$, get reward $r$, and reach $s'$.<br />
* To find out what action is best from $s'$, i.e. $a'$, we have to simulate all actions from $s'$. We can do this after we complete this iteration, then run $s',a_x$ for all $a_x\in\mathbb{A}$. But, we don't know how to complete this iteration without knowing this $a'$. So, another way is to simulate all actions from $s'$ using last known set of weights $\theta_{i-1}$. We just simulate state $s'$, action $a_x$ for all $a_x\in\mathbb{A}$ from the previous state and get $Q^*(s',a_x;\theta_{i-1})$. ('''Note''' that some papers do not use the set of weights from the previous iteration $\theta_{i-1}$. Instead they fix the weights for finding the best action for every $\tau$ steps to $\theta^-$, and do $Q^*(s',a_x;\theta^-)$ for $a_x\in\mathbb{A}$ and use this for the target value.)<br />
* Now we can compute our loss function using the Bellman equation, and backpropagate.<br />
$$<br />
\mbox{loss}=\mbox{target}-\mbox{prediction}=(r+\displaystyle\max_{a_x}Q^*(s',a_x;\theta_{i-1}))-Q^*(s,a;\theta_i)<br />
$$<br />
<br />
The '''problem''' with this approach is that at every iteration $i$, we have to do $|\mathbb{A}|$ forward passes on the previous set of weights $\theta_{i-1}$ to find out the best action $a'$ at $s'$. This becomes infeasible quickly with more possible actions.<br />
<br />
Authors of [[#References|[5]]] therefore use another kind of architecture. This architecture takes as input the state $s$, and computes the values $Q^*(s,a_x)$ for $a_x\in\mathbb{A}$. So there are $|\mathbb{A}|$ outputs. This basically parallelizes the forward passes so that $r+\displaystyle\max_{a_x}Q^*(s',a_x;\theta_{i-1})$ can be done with just a single pass through the outputs.<br />
<br />
<br />
[[File:DQN_arch.png|thumb||||600px|Source: [https://leonardoaraujosantos.gitbooks.io/artificial-inteligence/content/image_folder_7/DQNBreakoutBlocks.png leonardoaraujosantos.gitbooks.io]]]<br />
<br />
'''Note:''' When I say state $s$ as an input, I mean some representation of $s$. Since the environment is a partially observable MDP, it is hard to know $s$. So, we can for example, apply a CNN on the frames and get an idea of what the current state is. We pass this output to the input of the DNN (DNN is the fully connected layer for the CNN then).<br />
<br />
=== Experience Replay ===<br />
<br />
Authors of this paper borrow the concept of experience replay from [[#References|[5,6]]]. In experience replay, we do training in episodes. In each episode, we play and store consecutive $(s,a,r,s')$ tuples in the experience replay buffer. Then after the play, we choose random samples from this buffer and do our training.<br />
<br />
<br />
Advantages of experience replay over simple online Q learning<sup>[[#References|[5]]]</sup>:<br />
* '''Better data efficiency''': It is better to use one transition many times to learn again and again, rather than just learn once from it.<br />
* Learning from consecutive samples is difficult because of correlated data. Experience replay breaks this correlation.<br />
* Online learning means the input is decided by the previous action. So, if the maximising action is to go left in some game, next inputs would be about what happens when we go left. This can cause the optimiser to get stuck in a feedback loop, or even diverge, as [[#Reference|[7]]] points out.<br />
<br />
== Double Q Learning ==<br />
<br />
=== Problem with Q Learning<sup>[[#References|[8]]]</sup> ===<br />
<br />
For a simple neural network, each update tries to shift the current $Q^*$ estimate to a new value:<br />
<br />
$$<br />
Q^*(s,a) \leftarrow r+\gamma\displaystyle\max_{a_x}Q^*(s',a_x)<br />
$$<br />
<br />
Suppose the neural net has some inherent noise $\epsilon$. So, the neural net actually stores a value $\mathbb{Q}^*$ given by<br />
<br />
$$<br />
\mathbb{Q}^* = Q^*+\epsilon<br />
$$<br />
<br />
Even if $\epsilon$ has zero mean in the beginning, using the $\max$ operator at the update steps will start propagating $\gamma\cdot\max \mathbb{Q}^*$. This leads to a non zero mean subsequently. The problem is that "max causes overestimation because it does not preserve the zero-mean property of the errors of its operands." ([[#References|[8]]]) Thus, Q learning is more likely to choose overoptimistic values.<br />
<br />
=== How does Double Q Learning work? <sup>[[#References|[9]]]</sup> ===<br />
<br />
The problem can be solved by using two sets of weights $\theta$ and $\Theta$. The $\mbox{target}$ can be broken up as<br />
<br />
$$<br />
\mbox{target} = r+\displaystyle\max_{a_x}Q^*(s',a_x;\theta) = r+Q^*(s',\displaystyle\arg\max_{a_x}Q^*(s',a_x;\theta);\theta) = r+Q^*(s',a';\theta)<br />
$$<br />
<br />
Using double Q learning, we '''select''' the best action using current weights $\theta$ and '''evaluate''' the $Q^*$ value to decide the target value using $\Theta$.<br />
<br />
$$<br />
\mbox{target} = r+Q^*(s',\displaystyle\arg\max_{a_x}Q^*(s',a_x;\theta);\Theta) = r+Q^*(s',a';\Theta)<br />
$$<br />
<br />
This makes the evaluation fairer.<br />
<br />
=== Double Deep Q Learning ===<br />
<br />
[[#References|[9]]] further talks about how to use this for deep learning without much additional overhead. The suggestion is to use $\theta^-$ as $\Theta$.<br />
<br />
$$<br />
\mbox{target} = r+Q^*(s',\displaystyle\arg\max_{a_x}Q^*(s',a_x;\theta);\theta^-) = r+Q^*(s',a';\theta^-)<br />
$$<br />
<br />
== Bootstrapped DQN ==<br />
<br />
The authors propose an architecture that has a shared network and $K$ bootstrap heads. So, suppose our experience buffer $E$ has $n$ data points, where each datapoint is a $(s,a,r,s')$ tuple. Each bootstrap head trains on a different buffer $E_i$, where each $E_i$ has been constructed by sampling $n$ data points from the original experience buffer $E$ with replacement ('''bootstrap method''').<br />
<br />
<br />
Because each of the heads train on a different buffer, they model a different $Q^*$ function (say $Q^*_k$). Now, for each episode, we first choose a specific $Q^*_k=Q^*_s$. This $Q^*_s$ helps us create the experience buffer for the episode. From any state $s_t$, we populate the experience buffer by choosing the next action $a_t$ that maximises $Q^*_s$. (similar to '''Thompson Sampling''')<br />
<br />
$$<br />
a_t = \displaystyle\arg\max_a Q^*_s(s_t,a_t)<br />
$$<br />
<br />
Also, along with $s_t,a_t,r_t,s_{t+1}$, they push a bootstrap mask $m_t$. This mask is basically is a binary vector of size $K$, and it tells which $Q_k$ should be affected by this datapoint, if it is chosen as a training point. So, for example, if $K=5$ and there is a experience tuple $(s_t,a_t,r_t,s_{t+1},m_t)$ where $m_t=(0,1,1,0,1)$, then $(s_t,a_t,r_t,s_{t+1})$ should only affect $Q_2,Q_3$ and $Q_5$.<br />
<br />
<br />
So, at each iteration, we just choose few points from this buffer and train the respective $Q_{(\cdot)}$ based on the bootstrap masks.<br />
<br />
=== How to generate masks? ===<br />
<br />
Masks are created by sampling from the '''masking distribution'''. Now, there are many ways to choose this masking distribution:<br />
<br />
* If for each datapoint $D_i$ ($i=1$ to $n$), we mask from $\mbox{Bernoulli}(0.5)$, this will roughly allow us to have half the points from the original buffer. To get to size $n$, we duplicate these points by doubling the weights for each datapoint. This essentially gives us a '''double or nothing''' bootstrap<sup>[[#References|[10]]]</sup>.<br />
* If the mask is $(1, 1 \cdots 1)$, then this becomes an '''ensemble learning''' method.<br />
* $m_t~\mbox{Poi}(1)$ (poisson distribution)<br />
* $m_t[k]~\mbox{Exp}(1)$ (exponential distribution)<br />
<br />
For this paper's results, the authors used a $\mbox{Bernoulli}(p)$ distribution.<br />
<br />
== Related Work ==<br />
<br />
The authors mention the method described in [[#References|[11]]]. The authors of [[#References|[11]]] talk about the principle of "optimism in the face of uncertainty" and modify the reward function to encourage state-action pairs that have not been seen often:<br />
<br />
$$<br />
R(s,a) \leftarrow R(s,a)+\beta\cdot\mbox{novelty}(s,a)<br />
$$<br />
<br />
According to the authors, [[#References|[11]]]'s DQN algorithm relies on a lot of hand tuning and is only good for non stochastic problems. The authors further compare their results to [[#References|[11]]]'s results on Atari.<br />
<br />
<br />
The authors also mention an existing algorithm PSRL<sup>[[#References|[12,13]]]</sup>, or posterior sampling based RL. However, this algorithm requires a solved MDP, which is not feasible for large systems. Bootstrapped DQN approximates this idea by sampling from approximate $Q^*$ functions.<br />
<br />
<br />
Further, the authors mention that the work in [[#References|[12,13]]] has been followed by RLSVI<sup>[[#Reference|[14]]]</sup> which solves the problem for linear cases.<br />
<br />
== Deep Exploration: Why is Bootstrapped DQN so good at it? ==<br />
<br />
The authors consider a simple example to demonstrate the effectiveness of bootstrapped DQN at deep exploration.<br />
<br />
[[File:deep_exploration_example.png|thumb||center||700px|Source: this paper, section 5.1]]<br />
<br />
<br />
<br />
In this example, the agent starts at $s_2$. There are $N$ steps, and $N+9$ timesteps to generate the experience buffer. The agent is said to have learned the optimal policy if it achieves the best possible reward of $10$ (go to the rightmost state in $N-1$ timesteps, then stay there for $10$ timesteps), for at least $100$ such episodes. The results they got:<br />
<br />
[[File:deep_exploration_results.png|thumb||center||700px|Source: this paper, section 5.1]]<br />
<br />
<br />
<br />
The blue dots indicate when the agent learnt the optimal policy. If this took more than $2000$ episodes, they indicate it with a red dot. Thompson DQN is DQN with posterior sampling at every timestep. Ensemble DQN is same as bootstrapped DQN except that the mask is all $(1,1 \cdots 1)$. It is evident from the graphs that bootstrapped DQN can achieve deep exploration better than these two methods, and DQN.<br />
<br />
=== But why is it better? ===<br />
<br />
The authors say that this is because bootstrapped DQN constructs different approximations to the posterior $Q^*$ with the same initial data. This diversity of approximations is because of random initalization of weights for the $Q^*_k$ heads. This means that these heads start out trying random actions (because of diverse random initial $Q^*_k$), but when some head finds a good state and generalises to it, some (but not all) of the heads will learn from it, because of the bootstrapping. Eventually other heads will either find other good states, or end up learning the best good states found by the other heads.<br />
<br />
<br />
So, the architecture explores well and once a head achieves the optimal policy, eventually, all heads achieve the policy.<br />
<br />
== Results ==<br />
<br />
The authors test their architecture on 49 Atari games. They mention that there has been recent work to improve the performance of DDQNs, but those are tweaks whose intentions are orthogonal to this paper's idea. So, they don't compare their results with them.<br />
<br />
=== Scale: What values of $K$, $p$ are best? ===<br />
<br />
[[File:scale_k_p.png|thumb||center||800px|Source: this paper, section 6.1]]<br />
<br />
Recall that $K$ is the number of bootstrap heads and $p$ is the parameter for the masking distribution (Bernoulli). The authors say that around $K=10$, the performance reaches close to the peak, so it should be good.<br />
<br />
<br />
$p$ also represents the amount of data sharing. This is because lesser $p$ means there is lesser chance (due to the Bernoulli distribution) that the corresponding datapoint is taken into the bootstrapped dataset $D_i$. So, lesser $p$ means more identical datapoints, hence more heads share their datapoints. However, the value of $p$ doesn't seem to affect the rewards achieved over time. The authors give the following reasons for it:<br />
<br />
* The heads start with random weights for $Q^*$, so the targets (which use $Q^*$) turn out to be different. So the update rules are different.<br />
* Atari is deterministic.<br />
* Because of the initial diversity, the heads will learn differently even if they predict the same action for the given state.<br />
<br />
$p=1$ is the value they use finally, because this reduces the no. of identical datapoints and reduces time.<br />
<br />
=== Performance on Atari ===<br />
<br />
In general, the results tell us that bootstrapped DQN achieves better results.<br />
<br />
[[File:atari_results_bootstrapped_dqn.png|thumb||center||800px|Source: this paper, section 6.2]]<br />
<br />
The authors plot the improvement they achieved with bootstrapped DQN with the games. They define '''improvement''' to be $x$ if bootstrapped DQN achieves a better result than DQN in $\frac{1}{x}$ frames.<br />
<br />
[[File:bdqn_improvement.png|thumb||center||1000px|Source: this paper, section 6.2]]<br />
<br />
<br />
The authors say that bootstrapped DQN doesn't work good on all Atari games. They point out that there are some challenging games, where exploration is key but bootstrapped DQN doesn't do good enough (but does better than DQN). Some of these games are Frostbite and Montezuma’s Revenge. They say that even better exploration may help, but also point out that there may be other problems like: network instability, reward clipping and temporally extended rewards.<br />
<br />
=== Improvement: Highest Score Reached & how fast is this high score reached? ===<br />
<br />
The authors plot the improvement graphs after 20m and 200m frames.<br />
<br />
[[File:cumulative_rewards_bdqn.png|thumb||center||700px|Source: this paper, section 6.3]]<br />
<br />
=== Visualisation of Results ===<br />
<br />
One of the authors' [https://www.youtube.com/playlist?list=PLdy8eRAW78uLDPNo1jRv8jdTx7aup1ujM youtube playlist] can be found online.<br />
<br />
<br />
The authors also point out that just purely using bootstrapped DQN as an exploitative strategy is pretty good by itself, better than vanilla DQN. This is because of the deep exploration capabilities of bootstrapped DQN, since it can use the best states it knows and also plan to try out states it doesn't have any information about. Even in the videos, it can be seen that the heads agree at all the crucial decisions, but stay diverse at other less important steps.<br />
<br />
== Critique ==<br />
<br />
It would be very interesting and a great addition to the the experimental section of the paper, if the authors would have compared with asynchronous methods of exploration of the state space first introduced in [[#References|[15]]]. The authors unfortunately only compared their DQN with the original DQN and not all the other variations in the literature, and justified it by saying that their idea was "orthogonal" to these improvements.<br />
<br />
=== Different way to do exploration-exploitation? ===<br />
<br />
Instead of choosing the next action $a_t$ that maximises $Q^*_s$, they could have chosen different actions $a_i$ with probabilities<br />
<br />
$$<br />
\mathbb{P}(s_t,a_i) = \frac{Q^*_s(s_t,a_i)}{\displaystyle \sum_{i=1}^{|\mathbb{A}|} Q^*_s(s_t,a_i)}<br />
$$<br />
<br />
According to me, this is closer to Thompson Sampling.<br />
<br />
=== Why use Bernoulli? ===<br />
<br />
The choice of having a Bernoulli masking distribution eventually doesn't help them at all, since the algorithm does good because of the initial diversity. Maybe they can use some other masking distribution?<br />
<br />
=== Unanswered Questions & Miscellaneous ===<br />
* The Thompson DQN is not preferred because other randomized value functions can implement settings similar to Thompson sampling without the need for an intractable exact posterior update and also by working around the computational issue with Thompson Sampling: resampling every time step. Perhaps the authors could have explored Temporal Difference learning which is an attempt at combining Dynamic Programming and Monte Carlo methods.<br />
* The actual algorithm is hidden in the appendix. It could have been helpful if it were in the main paper.<br />
<br />
== References ==<br />
<br />
# [https://bandits.wikischolars.columbia.edu/file/view/Lecture+4.pdf Learning and optimization for sequential decision making, Columbia University, Lec 4]<br />
# [https://www.thoughtco.com/what-is-bootstrapping-in-statistics-3126172 Thoughtco, What is bootstrapping in statistics?]<br />
# [https://ocw.mit.edu/courses/mathematics/18-05-introduction-to-probability-and-statistics-spring-2014/readings/MIT18_05S14_Reading24.pdf Bootstrap confidence intervals, Class 24, 18.05, MIT Open Courseware]<br />
# [https://arxiv.org/abs/1506.02142 Yarin Gal and Zoubin Ghahramani. Dropout as a bayesian approximation: Representing model uncertainty in deep learning. arXiv preprint arXiv:1506.02142, 2015.]<br />
# [https://www.cs.toronto.edu/~vmnih/docs/dqn.pdf Mnih et al., Playing Atari with Deep Reinforcement Learning, 2015]<br />
# Long-Ji Lin. Reinforcement learning for robots using neural networks. Technical report, DTIC Document, 1993.<br />
# John N Tsitsiklis and Benjamin Van Roy. An analysis of temporal-difference learning with function approximation. Automatic Control, IEEE Transactions on, 42(5):674–690, 1997.<br />
# S. Thrun and A. Schwartz. Issues in using function approximation for reinforcement learning, 1993.<br />
# [https://arxiv.org/pdf/1509.06461.pdf Hado Van Hasselt, Arthur Guez, and David Silver. Deep reinforcement learning with double q-learning, 2015.]<br />
# [https://pdfs.semanticscholar.org/d623/c2cbf100d6963ba7dafe55158890d43c78b6.pdf Dean Eckles and Maurits Kaptein, Thompson Sampling with the Online Bootstrap, 2014, Pg 3]<br />
# [https://arxiv.org/abs/1507.00814 Bradly C. Stadie, Sergey Levine, Pieter Abbeel, Incentivizing Exploration In Reinforcement Learning With Deep Predictive Models, 2015.]<br />
# Ian Osband, Daniel Russo, and Benjamin Van Roy. (More) efficient reinforcement learning via posterior sampling, NIPS 2013.<br />
# Ian Osband and Benjamin Van Roy. Model-based reinforcement learning and the eluder dimension, NIPS 2014.<br />
# [https://arxiv.org/abs/1402.0635 Ian Osband, Benjamin Van Roy, Zheng Wen, Generalization and Exploration via Randomized Value Functions, 2014.]<br />
# Mnih, Volodymyr, et al. "Asynchronous methods for deep reinforcement learning." International Conference on Machine Learning. 2016.<br />
# George Konidaris, Sarah Osentoski, and Philip Thomas. 2011. Value function approximation in reinforcement learning using the fourier basis. In Proceedings of the Twenty-Fifth AAAI Conference on Artificial Intelligence (AAAI'11). AAAI Press 380-385.<br />
<br />
<br />
Other helpful links (unsorted):<br />
* [http://pemami4911.github.io/paper-summaries/deep-rl/2016/08/16/Deep-exploration.html pemami4911.github.io]<br />
* [http://www.stat.yale.edu/~pollard/Courses/241.fall97/Poisson.pdf Poisson Approximations]<br />
<br />
== Appendix ==<br />
<br />
=== Algorithm for Bootstrapped DQN ===<br />
The appendix lists the following algorithm. Periodically, the replay buffer is played back to update value function network Q.<br />
<br />
[[File:alg1.PNG|thumb||left||700px|Source: this paper's appendix]]</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Modular_Multitask_Reinforcement_Learning_with_Policy_Sketches&diff=30955Modular Multitask Reinforcement Learning with Policy Sketches2017-11-21T04:59:16Z<p>Nkadiyar: /* Introduction & Background */</p>
<hr />
<div>='''Introduction & Background'''=<br />
[[File:MRL0.png|border|right|400px]]<br />
[[File:MRL_diagram.jpg|thumb|right|400px| Figure 1b: the diagram for policy sketches]]<br />
[[File:MRL_encode.jpg|thumb|right|600px| Figure 1c: All sub tasks are encoded without any semantic meanings]]<br />
This paper describes a framework for learning composable deep subpolicies in a multitask setting. These policies are guided only by abstract sketches which are representative of the high-level behavior in the environment. Sketches annotate tasks with sequences of named subtasks, providing information about high-level structural relationships among tasks but not how to implement them—specifically not providing the detailed guidance used by much previous work on learning policy abstractions for RL (e.g. intermediate rewards, subtask completion signals, or intrinsic motivations). General reinforcement learning algorithms allow agents to solve tasks in complex environments. Vanilla policies find it difficult to deal with tasks featuring extremely delayed rewards. Most approaches often require in-depth supervision in the form of explicitly specified high-level actions, subgoals, or behavioral primitives. The proposed methodology is particularly suitable where rewards are difficult to engineer by hand. It is enough to tell the learner about the abstract policy structure, without indicating how high-level behaviors should try to use primitive percepts or actions.<br />
<br />
This paper explores a multitask reinforcement learning setting where the learner is presented with policy sketches. Policy sketches are defined as short, ungrounded, symbolic representations of a task. It describes its components, as shown in Figure 1. While symbols might be shared across different tasks ( the predicate "get wood" appears in sketches for both the tasks: "make planks" and "make sticks"). The learner is not shown or told anything about what these symbols mean, either in terms of observations or intermediate rewards. These sketches are initially meaningless but eventually the sketches get mapped into real policies. As shown in Figure 1c, the tasks are divided into human readable sub tasks. However, in the actual settings, the learner can only get access to encoded results.<br />
<br />
The agent learns from policy sketches by associating each high-level action with a parameterization of a low-level subpolicy. It jointly optimizes over concatenated task-specific policies by tying/sharing parameters across common subpolicies. They find that this architecture uses the high-level guidance provided by sketches to drastically accelerate learning of complex multi-stage behaviors. The experiments show that most benefits of learning from very detailed low-level supervision (e.g. from subgoal rewards) can also be obtained from fairly coarse high-level policy sketches. Most importantly, sketches are much easier to construct. They require no additions or modifications to the environment dynamics or reward function and can be easily provided by non-experts (third party mechanical turk providers). This makes it possible to extend the benefits of hierarchical RL to challenging environments where it may not be possible to specify by hand the details of relevant subtasks. This paper shows that their approach drastically outperforms purely unsupervised methods that do not provide the learner with any task-specific guidance. The specific use of sketches to parameterize modular subpolicies makes better use of sketches than conditioning on them directly.<br />
<br />
The modular structure of this whole approach, which associates every high-level action symbol with a discrete subpolicy, naturally leads to a library of interpretable policy fragments which can be are easily recombined. The authors evaluate the approach in a variety of different data conditions: <br />
# Learning the full collection of tasks jointly via reinforcement learning <br />
# In a zero-shot setting where a policy sketch is available for a held-out task<br />
# In an adaptation setting, where sketches are hidden and the agent must learn to use and adapt a pretrained policy to reuse high-level actions in a new task.<br />
<br />
The code has been released at http://github.com/jacobandreas/psketch.<br />
<br />
='''Related Work'''=<br />
The approach in this paper is a specific case of the options framework developed by Sutton et al., 1999. In that work, options are introduced as "closed-loop policies for taking action over the period of time". They show that options enable temporally abstract information to be included in reinforcement learning algorithms, though it was published before the large-scale popularity of neural networks for reinforcement.<br />
<br />
Other authors have recently explored techniques for learning policies which require less prior knowledge of the environment than the method presented in this paper. For example, in Vezhnevets et al. (2016), the authors propose a RNN architecture to build "implicit plans" only through interacting with the environment as in the classic reinforcement learning problem formulation.<br />
<br />
One closely related line of work is the Hierarchical Abstract Machines (HAM) framework introduced by Parr & Russell, 1998 [11]. Like the approach which the Modular Multitask Reinforcement Learning with Policy Sketches uses, HAMs begin with a representation of a high-level policy as an automaton (or a more general computer program; Andre & Russell,<br />
2001 [7]; Marthi et al., 2004 [12]) and use reinforcement learning to fill in low-level details.<br />
<br />
='''Learning Modular Policies from Sketches'''=<br />
The paper considers a multitask reinforcement learning problem arising from a family of infinite-horizon discounted Markov decision processes in a shared environment. This environment is specified by a tuple $(S, A, P, \gamma )$, with <br />
* $S$ a set of states<br />
* $A$ a set of low-level actions <br />
* $P : S \times A \times S \to R$ a transition probability distribution<br />
* $\gamma$ a discount factor<br />
<br />
Each task $t \in T$ is then specified by a pair $(R_t, \rho_t)$, with $R_t : S \to R$ a task-specific reward function and $\rho_t: S \to R$, an initial distribution over states. For a fixed sequence ${(s_i, a_i)}$ of states and actions obtained from a rollout of a given policy, we will denote the empirical return starting in state $s_i$ as $q_i = \sum_{j=i+1}^\infty \gamma^{j-i-1}R(s_j)$. In addition to the components of a standard multitask RL problem, we assume that tasks are annotated with sketches $K_t$ , each consisting of a sequence $(b_{t1},b_{t2},...)$ of high-level symbolic labels drawn from a fixed vocabulary $B$.<br />
<br />
==Model==<br />
The authors exploit the structural information provided by sketches by constructing for each symbol ''b'' a corresponding subpolicy $\pi_b$. By sharing each subpolicy across all tasks annotated with the corresponding symbol, their approach naturally learns the tied/shared abstraction for the corresponding subtask.<br />
<br />
[[File:Algorithm_MRL2.png|center|frame|Pseudo Algorithms for Modular Multitask Reinforcement Learning with Policy Sketches]]<br />
<br />
At every timestep, a subpolicy selects either a low-level action $a \in A$ or a special STOP action. The augmented state space is denoted as $A^+ := A \cup \{STOP\}$. At a high level, this framework is agnostic to the implementation of subpolicies: any function that takes a representation of the current state onto a distribution over $A^+$ will work fine with the approach.<br />
<br />
In this paper, $\pi_b$ is represented as a neural network. These subpolicies may be viewed as options of the kind described by [2], with the key distinction that they have no initiation semantics, but are instead invokable everywhere, and have no explicit representation as a function from an initial state to a distribution over final states (instead this paper uses the STOP action to terminate).<br />
<br />
Given a fixed sketch $(b_1, b_2,....)$, a task-specific policy $\Pi_r$ is formed by concatenating its associated subpolicies in sequence. In particular, the high-level policy maintains a sub-policy index ''i'' (initially 0), and executes actions from $\pi_{b_i}$ until the STOP symbol is emitted, at which point control is passed to bi+1 . We may thus think of as inducing a Markov chain over the state space $S \times B$, with transitions:<br />
[[File:MRL1.png|center|border|]]<br />
<br />
Note that $\Pi_r$ is semi-Markov with respect to projection of the augmented state space $S \times B$ onto the underlying state space ''S''. The complete family of task-specific policies is denoted as $\Pi := \bigcup_r \{ \Pi_r \}$. Assume each $\pi_b$ be an arbitrary function of the current environment state parameterized by some weight vector $\theta_b$. The learning problem is to optimize overall $\theta_b$ to maximize expected discounted reward<br />
[[File:MRL2.png|center|border|]]<br />
across all tasks $t \in T$.<br />
<br />
==Policy Optimization==<br />
<br />
Control policy parameterized by parameter vector $\theta$, $$\displaystyle \max_{\Theta}E[\sum_{t=0}^{H}R(s_{t})|\pi_{\theta}]$$ $\pi_{\theta}(u|s)$ is the probability of action u in state s. The details of policy optimization can be found here: https://people.eecs.berkeley.edu/~pabbeel/nips-tutorial-policy-optimization-Schulman-Abbeel.pdf<br />
<br />
Here that optimization is accomplished through a simple decoupled actor–critic method. In a standard policy gradient approach, with a single policy $\pi$ with parameters $\theta$, the gradient steps are of the form:<br />
[[File:MRL3.png|center|border|]]<br />
<br />
where the baseline or “critic” c can be chosen independently of the future without introducing bias into the gradient. Recalling the previous definition of $q_i$ as the empirical return starting from $s_i$, this form of the gradient corresponds to a generalized advantage estimator with $\lambda = 1$. Here ''c'' achieves close to the optimal variance[6] when it is set exactly equal to the state-value function $V_{\pi} (s_i) = E_{\pi} q_i$ for the target policy $\pi$ starting in state $s_i$.<br />
[[File:MRL4.png|frame|]]<br />
<br />
In the case of generalizing to modular policies built by sequencing sub-policies the authors suggest to have one subpolicy per symbol but one critic per task. This is because subpolicies $\pi_b$ might participate in many compound policies $\Pi_r$, each associated with its own reward function $R_r$ . Thus individual subpolicies are not uniquely identified or differentiated with value functions. The actor–critic method is extended to allow decoupling of policies from value functions by allowing the critic to vary per-sample (per-task-and-timestep) based on the reward function with which that particular sample is associated. Noting that <br />
[[File:MRL5.png|center|border|]]<br />
i.e. the sum of gradients of expected rewards across all tasks in which $\pi_b$ participates, we have:<br />
[[File:MRL6.png|center|border|]]<br />
where each state-action pair $(s_{t_i}, a_{t_i})$ was selected by the subpolicy $\pi_b$ in the context of the task ''t''.<br />
<br />
Now minimization of the gradient variance requires that each $c_t$ actually depend on the task identity. (This follows immediately by applying the corresponding argument in [6] individually to each term in the sum over ''t'' in Equation 2.) Because the value function is itself unknown, an approximation must be estimated from data. Here it is allowed that these $c_t$ to be implemented with an arbitrary function approximator with parameters $\eta_t$ . This is trained to minimize a squared error criterion, with gradients given by<br />
[[File:MRL7.png|center|border|]]<br />
Alternative forms of the advantage estimator (e.g. the TD residual $R_t (s_i) + \gamma V_t(s_{i+1} - V_t(s_i))$ or any other member of the generalized advantage estimator family) can be used to substitute by simply maintaining one such estimator per task. Experiments show that conditioning on both the state and the task identity results in dramatic performance improvements, suggesting that the variance reduction given by this objective is important for efficient joint learning of modular policies.<br />
<br />
The complete algorithm for computing a single gradient step is given in Algorithm 1. (The outer training loop over these steps, which is driven by a curriculum learning procedure, is shown in Algorithm 2.) Note that this is an on-policy algorithm. In every step, the agent samples tasks from a task distribution provided by a curriculum (described in the following subsection). The current family of policies '''$\Pi$''' is used to perform rollouts for every sampled task, accumulating the resulting tuples of (states, low-level actions, high-level symbols, rewards, and task identities) into a dataset ''$D$''. Once ''$D$'' reaches a maximum size D, it is used to compute gradients with respect to both policy and critic parameters, and the parameter vectors are updated accordingly. The step sizes $\alpha$ and $\beta$ in Algorithm 1 can be chosen adaptively using any first-order method.<br />
<br />
==Curriculum Learning==<br />
<br />
For complex tasks, like the one depicted in Figure 3b, it is difficult for the agent to discover any states with positive reward until many subpolicy behaviors have already been learned. It is thus a better use of the learner’s time (and computational resources) to focus on “easy” tasks, where many rollouts will result in high reward from which relevant subpolicy behavior can be obtained. But there is a fundamental tradeoff involved here: if the learner spends a lot of its time on easy tasks before being told of the existence of harder ones, it may overfit and learn subpolicies that exhibit the desired structural properties or no longer generalize.<br />
<br />
To resolve these issues, a curriculum learning scheme is used that allows the model to smoothly scale up from easy tasks to more difficult ones without overfitting. Initially, the model is presented with tasks associated with short sketches. Once average reward on all these tasks reaches a certain threshold, the length limit is incremented. It is assumed that rewards across tasks are normalized with maximum achievable reward $0 < q_i < 1$ . Let $Er_t$ denote the empirical estimate of the expected reward for the current policy on task T. Then at each timestep, tasks are sampled in proportion $1-Er_t$, which by assumption must be positive.<br />
<br />
Intuitively, the tasks that provide the strongest learning signal are those in which <br />
# The agent does not on average achieve reward close to the upper bound<br />
# Many episodes result in a high reward.<br />
<br />
The expected reward component of the curriculum solves condition (1) by making sure that time is not spent on nearly solved tasks, while the length bound component of the curriculum addresses condition (2) by ensuring that tasks are not attempted until high-reward episodes are likely to be encountered. The experiments performed show that both components of this curriculum learning scheme improve the rate at which the model converges to a good policy.<br />
<br />
The complete curriculum-based training algorithm is written as Algorithm 2 above. Initially, the maximum sketch length $l_{max}$ is set to 1, and the curriculum initialized to sample length-1 tasks uniformly. For each setting of $l_{max}$, the algorithm uses the current collection of task policies to compute and apply the gradient step described in Algorithm 1. The rollouts obtained from the call to TRAIN-STEP can also be used to compute reward estimates $Er_t$ ; these estimates determine a new task distribution for the curriculum. The inner loop is repeated until the reward threshold $r_{good}$ is exceeded, at which point $l_{max}$ is incremented and the process repeated over a (now-expanded) collection of tasks.<br />
<br />
='''Experiments'''=<br />
[[File:MRL8.png|border|right|400px]]<br />
This paper considers three families of tasks: a 2-D Minecraft-inspired crafting game (Figure 3a), in which the agent must acquire particular resources by finding raw ingredients, combining them together in the correct order, and in some cases building intermediate tools that enable the agent to alter the environment itself; a 2-D maze navigation task that requires the agent to collect keys and open doors, and a 3-D locomotion task (Figure 3b) in which a quadrupedal robot must actuate its joints to traverse a narrow winding cliff.<br />
<br />
In all tasks, the agent receives a reward only after the final goal is accomplished. For the most challenging tasks, involving sequences of four or five high-level actions, a task-specific agent initially following a random policy essentially never discovers the reward signal, so these tasks cannot be solved without considering their hierarchical structure. These environments involve various kinds of challenging low-level control: agents must learn to avoid obstacles, interact with various kinds of objects, and relate fine-grained joint activation to high-level locomotion goals.<br />
<br />
==Implementation==<br />
In all of the experiments, each subpolicy is implemented as a neural network with ReLU nonlinearities and a hidden layer with 128 hidden units. Each critic is a linear function of the current state. Each subpolicy network receives as input a set of features describing the current state of the environment and outputs a distribution over actions. The agent acts at every timestep by sampling from this distribution. The gradient steps given in lines 8 and 9 of Algorithm 1 are implemented using RMSPROP with a step size of 0.001 and gradient clipping to a unit norm. They take the batch size D in Algorithm 1 to be 2000, and set $\gamma$= 0.9 in both environments. For curriculum learning, the improvement threshold $r_{good}$ is 0.8.<br />
<br />
==Environments==<br />
<br />
The environment in Figure 3a is inspired by the popular game Minecraft, but is implemented in a discrete 2-D world. The agent interacts with objects in the environment by executing a special USE action when it faces them. Picking up raw materials initially scattered randomly around the environment adds to an inventory. Interacting with different crafting stations causes objects in the agent’s inventory to be combined or transformed. Each task in this game corresponds to some crafted object the agent must produce; the most complicated goals require the agent to also craft intermediate ingredients, and in some cases build tools (like a pickaxe and a bridge) to reach ingredients located in initially inaccessible regions of the world.<br />
<br />
[[File:MRL_maze.png|boarder|right|400px]]<br />
<br />
The maze environment is very similar to “light world” described by [4], which can be seen in Figure 3c. The agent is placed in a discrete world consisting of a series of rooms, some of which are connected by doors. The agent needs to first pick up a key to open them. For our experiments, each task corresponds to a goal room that the agent must reach through a sequence of intermediate rooms. The agent senses the distance to keys, closed doors, and open doors in each direction. Sketches specify a particular sequence of directions for the agent to traverse between rooms to reach the goal. The sketch always corresponds to a viable traversal from the start to the goal position, but other (possibly shorter) traversals may also exist.<br />
<br />
The cliff environment (Figure 3b) proves the effectiveness of the approach in a high-dimensional continuous control environment where a quadrupedal robot [5] is placed on a variable-length winding path, and must navigate to the end without falling off. This is a challenging RL problem since the walker must learn the low-level walking skill before it can make any progress. The agent receives a small reward for making progress toward the goal, and a large positive reward for reaching the goal square, with a negative reward for falling off the path.<br />
<br />
==Multitask Learning==<br />
<br />
[[File:MRL9.png|border|center|800px]]<br />
The primary experimental question in this paper is whether the extra structure provided by policy sketches alone is enough to enable fast learning of coupled policies across tasks. The aim is to explore the differences between the approach described and relevant prior work that performs either unsupervised or weakly supervised multitask learning of hierarchical policy structure. Specifically, they compare their '''modular''' approach to:<br />
<br />
# Structured hierarchical reinforcement learners:<br />
#* the fully unsupervised '''option–critic''' algorithm of Bacon & Precup[1]<br />
#* a '''Q automaton''' that attempts to explicitly represent the Q function for each task/subtask combination (essentially a HAM [8] with a deep state abstraction function)<br />
# Alternative ways of incorporating sketch data into standard policy gradient methods:<br />
#* learning an '''independent''' policy for each task<br />
#* learning a '''joint policy''' across all tasks, conditioning directly on both environment features and a representation of the complete sketch<br />
<br />
The joint and independent models performed best when trained with the same curriculum described in Section 3.3, while the option–critic model performed best with a length–weighted curriculum that has access to all tasks from the beginning of training.<br />
<br />
Learning curves for baselines and the modular model is shown in Figure 4. It can be seen that in all environments, our approach substantially outperforms the baselines: it induces policies with substantially higher average reward and converges more quickly than the policy gradient baselines. It can further be seen in Figure 4c that after policies have been learned on simple tasks, the model is able to rapidly adapt to more complex ones, even when the longer tasks involve high-level actions not required for any of the short tasks.<br />
<br />
==Ablations==<br />
[[File:MRL10.png|border|right|400px]]<br />
In addition to the overall modular parameter tying structure induced by sketches, the other critical component of the training procedure is the decoupled critic and the curriculum. The next experiments investigate the extent to which these are necessary for good performance.<br />
<br />
To evaluate the critic, consider three ablations: <br />
# Removing the dependence of the model on the environment state, in which case the baseline is a single scalar per task<br />
# Removing the dependence of the model on the task, in which case the baseline is a conventional generalized advantage estimator<br />
# Removing both, in which case the baseline is a single scalar, as in a vanilla policy gradient approach.<br />
<br />
Results are shown in Figure 5a. Introducing both state and task dependence into the baseline leads to faster convergence of the model: the approach with a constant baseline achieves less than half the overall performance of the full critic after 3 million episodes. Introducing task and state dependence independently improve this performance; combining them gives the best result.<br />
<br />
Two other experiments are also performed as Figure 5b: starting with short examples and moving to long ones, and sampling tasks in inverse proportion to their accumulated reward. It is shown that both components help; prioritization by both length and weight gives the best results.<br />
<br />
==Zero-shot and Adaptation Learning==<br />
[[File:MRL11.png|border|left|320px]]<br />
In the final experiments, the authors test the model’s ability to generalize beyond the standard training condition. Consider two tests of generalization: a zero-shot setting, in which the model is provided a sketch for the new task and must immediately achieve good performance, and an adaptation setting, in which no sketch is provided leaving the model to learn the form of a suitable sketch via interaction in the new task.They hold out two length-four tasks from the full inventory used in Section 4.3, and train on the remaining tasks. For zero-shot experiments, the concatenated policy is formed to describe the sketches of the held-out tasks, and repeatedly executing this policy (without learning) in order to obtain an estimate of its effectiveness. For adaptation experiments, consider ordinary RL over high-level actions B rather than low-level actions A, implementing the high-level learner with the same agent architecture as described in Section 3.1. Results are shown in Table 1. The held-out tasks are sufficiently challenging that the baselines are unable to obtain more than negligible reward: in particular, the joint model overfits to the training tasks and cannot generalize to new sketches, while the independent model cannot discover enough of a reward signal to learn in the adaptation setting. The modular model does comparatively well: individual subpolicies succeed in novel zero-shot configurations (suggesting that they have in fact discovered the behavior suggested by the semantics of the sketch) and provide a suitable basis for adaptive discovery of new high-level policies.<br />
<br />
='''Conclusion & Critique'''=<br />
The paper's contributions are:<br />
<br />
* A general paradigm for multitask, hierarchical, deep reinforcement learning guided by abstract sketches of task-specific policies.<br />
<br />
* A concrete recipe for learning from these sketches, built on a general family of modular deep policy representations and a multitask actor–critic training objective.<br />
<br />
They have described an approach for multitask learning of deep multitask policies guided by symbolic policy sketches. By associating each symbol appearing in a sketch with a modular neural sub policy, they have shown that it is possible to build agents that share behavior across tasks in order to achieve success in tasks with sparse and delayed rewards. This process induces an inventory of reusable and interpretable sub policies which can be employed for zero-shot generalization when further sketches are available, and hierarchical reinforcement learning when they are not.<br />
<br />
Hierarchical Reinforcement Learning is a popular research topic now, it is interesting to compare this paper(which was presented in this class)[13], where the architecture follows the manager-and-workers style. In that work, the sub policy is decided by manager network. To finish a hierarchical task, each worker focuses on the sub task and optimizes the sub task. The difference is in that work, the sub policy is implicit and is trained in the training process. A further question for future work for this paper: would these sub policy be learned automatically instead of pre-defined?<br />
<br />
There are four drawbacks of the presented work. First, the presented ideas portrayed in this paper (for instance: symbolic specifications, actor-critic, shared representations) have all been explored in other works. Secondly, aforementioned approach relies heavily on curriculum learning - which is, as we know, difficult to design as it is pretty complicated depending on the task in-hand. Thirdly, there has been no discussion on how curriculum can be designed for larger problems. Finally, this approach could be that building of different neural networks for each sub tasks could lead to overly complicated networks and is not in the spirit of building an efficient structure.<br />
<br />
= '''Resources''' =<br />
You can find a talk on this paper [https://www.google.ca/url?sa=t&rct=j&q=&esrc=s&source=video&cd=3&cad=rja&uact=8&ved=0ahUKEwjNzPuBqM7XAhVK6mMKHQICAdEQtwIILDAC&url=https%3A%2F%2Fwww.youtube.com%2Fwatch%3Fv%3DNRIcDEB64x8&usg=AOvVaw1NHi2XExGXwhzzeJn5AcnR here].<br />
<br />
<br />
<br />
='''References'''=<br />
[1] Bacon, Pierre-Luc and Precup, Doina. The option-critic architecture. In NIPS Deep Reinforcement Learning Work-shop, 2015.<br />
<br />
[2] Sutton, Richard S, Precup, Doina, and Singh, Satinder. Be-tween MDPs and semi-MDPs: A framework for temporal abstraction in reinforcement learning. Artificial intel-ligence, 112(1):181–211, 1999.<br />
<br />
[3] Stolle, Martin and Precup, Doina. Learning options in reinforcement learning. In International Symposium on Abstraction, Reformulation, and Approximation, pp. 212– 223. Springer, 2002.<br />
<br />
[4] Konidaris, George and Barto, Andrew G. Building portable options: Skill transfer in reinforcement learning. In IJ-CAI, volume 7, pp. 895–900, 2007.<br />
<br />
[5] Schulman, John, Moritz, Philipp, Levine, Sergey, Jordan, Michael, and Abbeel, Pieter. Trust region policy optimization. In International Conference on Machine Learning, 2015b.<br />
<br />
[6] Greensmith, Evan, Bartlett, Peter L, and Baxter, Jonathan. Variance reduction techniques for gradient estimates in reinforcement learning. Journal of Machine Learning Research, 5(Nov):1471–1530, 2004.<br />
<br />
[7] Andre, David and Russell, Stuart. Programmable reinforce-ment learning agents. In Advances in Neural Information Processing Systems, 2001.<br />
<br />
[8] Andre, David and Russell, Stuart. State abstraction for pro-grammable reinforcement learning agents. In Proceedings of the Meeting of the Association for the Advance-ment of Artificial Intelligence, 2002.<br />
<br />
[9] Author Jacob Andreas presenting the paper - https://www.youtube.com/watch?v=NRIcDEB64x8<br />
<br />
[10] Vezhnevets, A., Mnih, V., Osindero, S., Graves, A., Vinyals, O., & Agapiou, J. (2016). Strategic attentive writer for learning macro-actions. In Advances in Neural Information Processing Systems (pp. 3486-3494).<br />
<br />
[11] Parr, Ron and Russell, Stuart. Reinforcement learning with hierarchies of machines. In Advances in Neural Information Processing Systems, 1998.<br />
<br />
[12] Marthi, Bhaskara, Lantham, David, Guestrin, Carlos, and Russell, Stuart. Concurrent hierarchical reinforcement learning. In Proceedings of the Meeting of the Association for the Advancement of Artificial Intelligence, 2004.<br />
<br />
[13] A. S. Vezhnevets, S. Osindero, T. Schaul, N. Heess, M. Jaderberg, D. Silver, and K. Kavukcuoglu. Feudal networks for hierarchical reinforcement learning. arXiv preprint arXiv:1703.01161, 2017.<br />
<br />
= Appendix =<br />
The authors provide a brief appendix that gives a complete list of tasks and sketches. Asterisk * indicates that the task was held out for generalization experiments in Section 4.5, but included in the multitask experiments of Sections 4.3 and 4.4.<br />
<br />
[[File: tasks.PNG]]</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Universal_Style_Transfer_via_Feature_Transforms&diff=30831Universal Style Transfer via Feature Transforms2017-11-19T19:20:50Z<p>Nkadiyar: /* Content/Style Balance */</p>
<hr />
<div>=Introduction=<br />
When viewing an image, whether it is a photograph or a painting, two types of mutually exclusive data are present. First, there is the content of the image, such as a person in a portrait. However, the content does not uniquely define the image. Consider a case where multiple artists paint a portrait of an identical subject, the results would vary despite the content being invariant. The cause of the variance is rooted in the style of each particular artist. Therefore, style transfer between two images results in the content being unaffected but the style being copied. Style transfer is an important image editing task which enables the creation of new artistic works. Typically one image is termed the content/reference image, whose style is discarded. The other image is called the style image, whose style, but the not content is copied to the content image.<br />
<br />
Deep learning techniques have been shown to be effective methods for implementing style transfer. Previous methods have been successful but with several key limitations and often trade off between generalization, quality and efficiency. Either they are fast, but have very few styles that can be transferred or they can handle arbitrary styles but are no longer efficient. The presented paper establishes a compromise between these two extremes by using only whitening and coloring transforms (WCT) to transfer a style within a feedforward image reconstruction architecture. No training of the underlying deep network is required per style.<br />
<br />
=Related Work=<br />
Gatys et al. developed a new method for generating textures from sample images in 2015 [1] and extended their approach to style transfer by 2016 [2]. They proposed the use of a pre-trained convolutional neural network (CNN) to separate content and style of input images. Having proven successful, a number of improvements quickly developed, reducing computational time, increasing the diversity of transferrable styes, and improving the quality of the results. Central to these approaches and of the present paper is the use of a CNN.<br />
<br />
In 2017, Mechrez et al. [12] proposed an approach that takes as input a stylized image and makes it more photorealistic. Their approach relied on the Screened Poisson Equation, maintaining the fidelity of the stylized image while constraining the gradients to those of the original input image. The method they proposed was fast, simple, fully automatic and showed positive progress in making a stylized image photorealistic.<br />
<br />
Alternative attempts, by using a single network to transfer<br />
multiple styles, include models conditioned on binary selection units [13], a network that learns a set of new filters for every new style [15], and a novel conditional normalization layer that learns normalization parameters for each style [3]<br />
==How Content and Style are Extracted using CNNs==<br />
A CNN was chosen due to its ability to extract high level feature from images. These features can be interpreted in two ways. Within layer <math> l </math> there are <math> N_l </math> feature maps of size <math> M_l </math>. With a particular input image, the feature maps are given by <math> F_{i,j}^l </math> where <math> i </math> and <math> j </math> locate the map within the layer. Starting with a white noise image and an reference (content) image, the features can be transferred by minimizing<br />
<br />
<center><br />
<math> \mathcal{L}_{content} = \frac{1}{2} \sum_{i,j} \left( F_{i,j}^l - P_{i,j}^l \right)^2 </math><br />
</center><br />
<br />
where <math> P_{i,j} </math> denotes the feature map output caused by the white noise image. Therefore this loss function preserves the content of the reference image. The style is described using a Gram matrix given by<br />
<br />
<center><br />
<math><br />
G_{i,j}^l = \sum_k F_{i,k}^l F_{j,k}^l<br />
</math><br />
</center><br />
<br />
Gram matrix $G$ of a set of vectors $v_1,\dots,v_n$ is the matrix of all possible inner products whose entries are given by $G_{ij}=v_i^Tv_j$. The loss function that describes a difference in style between two images is equal to:<br />
<br />
<center><br />
<math><br />
\mathcal{L}_{style} = \frac{1}{4 N_l^2 M_l^2} \sum_{i,j} \left(G_{i,j}^l - A_{i,j}^l \right)^2<br />
</math><br />
</center><br />
<br />
where <math> A_{i,j}^l </math> and <math> G_{i,j}^l </math> are the Gram matrices of the generated image and style image respectively. Therefore three images are required, a style image, a content image and an initial white noise image. Iterative optimization is then used to add content from one image to the white noise image, and style from the other. An additional parameter is used to balance the ratio of these loss functions.<br />
<br />
The 19-layer ImageNet trained VGG network was chosen by Gatys et al. VGG-19 is still commonly used in more recent works as will be shown in the presented paper, although training datasets vary. Such CNNs are typically used in classification problems by finalizing their output through a series of full connected layers. For content and style extraction it is the convolutional layers that are required. The method of Gatys et al. is style independent, since the CNN does not need to be trained for each style image. However the process of iterative optimization to generate the output image is computational expensive.<br />
<br />
==Other Methods==<br />
Other methods avoid the inefficiency of iterative optimization by training a network/networks on a set of styles. The network then directly transfers the style from the style image to the content image without solving the iterative optimization problem. V. Dumoulin et al. trained a single network on $N$ styles [3]. This improved upon previous work where a network was required per style [4]. The stylized output image was generated by simply running a feedforward pass of the network on the content image. While efficiency is high, the method is no longer able to apply an arbitrary style without retraining.<br />
<br />
=Methodology=<br />
Li et al. have proposed a novel method for generating the stylized image. A CNN is still used as in Gatys et al. to extract content and style. However, the stylized image is not generated through iterative optimization or a feed-forward pass as required by previous methods. Instead, whitening and colour transforms are used.<br />
<br />
==Image Reconstruction==<br />
[[File:image_resconstruction.png|thumb|150px|right|alt=Training a single decoder.|Training a single decoder. X denotes the layer of the VGG encoder that the decoder receives as input.]]<br />
An auto-encoder network is used to first encode an input image into a set of feature maps, and then decode it back to an image as shown in the adjacent figure. The encoder network used is VGG-19. This network is reponsible for obtaining feature maps (similar to Gatys et al.). The output of each of the first five layers is then fed into a corresponding decoder network, which is a mirrored version of VGG-19. Each decoder network then decodes the feature maps of the $l$th layer producing an output image. A mechanism for transferring style will be implemented by manipulating the feature maps between the encoder and decoder networks.<br />
<br />
First, the auto-encoder network needs to be trained. The following loss function is used<br />
<br />
<center><br />
<math><br />
<br />
\mathcal{L} = || I_{output} - I_{input} ||_2^2 + \lambda || \Phi(I_{output}) - \Phi(I_{input})||_2^2<br />
<br />
</math><br />
</center><br />
<br />
where $I_{input}$ and $I_{output}$ are the input and output images of the auto-encoder. $\Phi$ is the VGG encoder. The first term of the loss is the pixel reconstruction loss, while the second term is feature loss. Recall from "Related Work" that the feature maps correspond to the content of the image. Therefore the second term can also be seen as penalising for content differences that arise due to the encoder network. The network was trained using the Microsoft COCO dataset. <br />
<br />
They use whitening and coloring transforms to directly transform the $f_c$ (VGG feature<br />
map of content image at a certain layer) to match the covariance matrix of $f_s$ (VGG feature<br />
map of style image). This process is consisted of two steps, i.e., whitening and coloring transform. Note that the decoder will reconstruct the original content image if $f_c$ is directly fed into it.<br />
<br />
==Whitening Transform==<br />
Whitening first requires that the covariance of the data is a diagonal matrix. This is done by solving for the covariance matrix's eigenvalues and eigenvector matrices. Whitening then forces the diagonal elements of the eigenvalue matrix to be the same. This is achieved for a feature map from VGG through the following steps.<br />
<br />
# The feature map $f_c$ is extracted from a layer of the encoder network after activation on the content image. This is the data to be whitened.<br />
# $f_c$ is centered by subtracting its mean vector $m_c$.<br />
# Then, the eigenvectors $E_c$ and eigenvalues $D_c$ are found for the covariance matrix of $f_c$. <br />
# The whitened feature map is then given by $\hat{f}_c = E_c D_c^{-1/2} E_c^T f_c$. <br />
<br />
If interested, the derivation of the whitening equation can be seen in [5]. Li et al. found that whitening removed styles from the image.<br />
<br />
==Colour Transform==<br />
However, whitening does not transfer style from the style image. It only uses feature maps from the content image. The colour transform uses both $\hat{f}_c$ from above and $f_s$, the feature map from the style image.<br />
<br />
# $f_s$ is centered by subtracting its mean vector $m_s$.<br />
# Then, the eigenvectors $E_s$ and eigenvalues $D_s$ are calculated for the covariance matrix of $f_s$.<br />
# The colour transform is given by $\hat{f}_{cs} = E_s D_s^{1/2} E_s^T \hat{f}_c$.<br />
# Recenter $\hat{f}_{cs}$ using $m_s$.<br />
<br />
Intuitively, colouring results in a correlation between the $\hat{f}_c$ and $f_s$ feature maps. This is where the style transfer takes place.<br />
<br />
==Content/Style Balance==<br />
Using just $\hat{f}_{cs}$ as the input to the decoder may create a result that is too extreme in style. To balance content and style a new parameter $\alpha$ is defined.<br />
<br />
<center><br />
<math><br />
<br />
\hat{f}_{cs} = \alpha \hat{f}_{cs} + (1 - \alpha) f_c<br />
<br />
</math><br />
</center><br />
<br />
Authors use $\alpha$ = 0.6 in the style transfer experiments.<br />
<br />
==Using Multiple Layers==<br />
It has been previously mentioned that multiple decoders were trained, one for each of the first five layers of the encoder network. Each layer of a CNN perceives features at different levels. Levels close to the input image will detect lower level local features such as edges. Those levels deeper into the network will detect more complex global features. The style transfer algorithm is applied at each of these levels, which yields the question as to which results, as shown below, to use.<br />
<br />
[[File:multilevel_features.png|thumb|700px|center|alt=Results of style transfer from each of the first five layers of the encoder network.|Results of style transfer from each of the first five layers of the encoder network.]]<br />
<br />
Ideally, the results of each layer should be used to build the final output image. This captures the entire range of features detected by the encoder network. First, one full pass of the network is performed. Then the stylised image from the deepest layer (Relu_5_1 in this case) is taken and used as the content image for another iteration of the algorithm, where then the next layer (Relu_4_1) is used as the output. These steps are repeated until the final image is produced from the shallowest layer. This process is summarised in the figure below.<br />
<br />
[[File:process_summary.png|thumb|700px|center|alt=Process summary of the multi-level stylization algorithm.|The content (C) and style (S) are fed to the VGG encoding network. The output image (I) after a whitening and colour transform (WCT) is taken from the deepest level's decoder. The process is iteratively repeated until the most shallow layer is reached.]]<br />
<br />
The authors note that the transformations must be applied first at the highest level (most abstract) layers, which capture complicated local structures and pass this transformed image to lower layers, which improve on details. They observe that reversing this order (lowest to highest) leads to images with low visual quality, as low-level information cannot be preserved after manipulating high level features.<br />
<br />
[[File:Universal_Style_Transfer_Coarse_to_Fine.JPG|thumb|700px|center|alt=(a)-(c) Output from intermediate layers. (d) Reversed transformation order.|(a)-(c) Output from intermediate layers. (d) Reversed transformation order.]]<br />
<br />
=Evaluation=<br />
The success of style transfer might appear hard to quantify as it relies on qualitative judgement. However, the extremes of transferring no style, or transferring only style can be considered as performing poorly. Consistent transfer of style throughout the entire image is another parameter of success. Ideally, the viewer can recognize the content of the image, while seeing it expressed in an alternative style. Quantitatively, the quality of the style transfer can be calculated by taking the covariance matrix difference $L_s$ between the resulting image and the original style. The results of the presented paper also need to be considered within the contexts of generality, efficiency and training requirements.<br />
<br />
==Style Transfer==<br />
A number of style transfer examples are presented relative to other works. <br />
<br />
[[File:transfer_results_label.jpg|thumb|700px|center|alt=Style transfer results of the presented paper.|A: See [6]. B: See [7]. C: See [8]. D: Gatys et al. iterative optimization, see [2]. E: This paper's results.]]<br />
<br />
Li et al. then obtained the average $L_s$ using 10 random content images across 40 style images. They had the lowest average $log(L_s)$ of all referenced works at 6.3. Next lowest was Gatys et al. [2] with $log(L_s) = 6.7$. It should be noted that while $L_s$ quantitatively calculates the success of the style transfer, results are still subject to the viewer's impression. Reviewing the transfer results, rows five and six for Gatys et al.'s method shows local minimization issues. However, their method still achieves a competitive $L_s$ score.<br />
<br />
==Transfer Efficiency==<br />
It was hypothesized by Li et al. that using WCT would enable faster run-times than [2] while still supporting arbitrary style transfer. For a 256x256 image, using a 12GB TITAN X, they achieved a transfer time of 1.5 seconds. Gatys et al.'s method [2] required 21.2 seconds. The pure feed-forward approaches [7], and [8] had times equal to or less than 0.2 seconds. [6] had a time comparable to the presented paper's method. However, [6,7,8] do not generalize well to multiple styles as training is required. Therefore this paper obtained a near 15x speed up for a style agnostic transfer algorithm when compared to leading previous work. The authors also note that WCT was done using the CPU. They intend to port WCT to the GPU and expect to see the computational time be further reduced.<br />
<br />
==Other Applications==<br />
Li et al.'s method can also be used for texture synthesis. This was the original work of Gatys et. al. before they applied their algorithm to style transfer problems. Texture synthesis takes a reference texture/image and creates new textures from it. With proper boundary conditions enforced these synthesized textures can be tileable. Alternatively, higher resolution textures can be generated. Texture synthesis has applications in areas such as computer graphics, allowing for large surfaces to be texture mapped.<br />
<br />
The content image is set as white noise, similar to how [2] initializes their output image. Then the reference texture/image is set as the style image. Since the content image is initially random white noise, then the features generated by the encoder of this image are also random. Li et al. state that this increases the diversity of the resulting output textures.<br />
<br />
[[File:texture_synthesis_label.jpg|thumb|700px|center|alt=Texture synthesis results.|A: Reference image/texture. B: Result from [8]. C: Result of present paper.]]<br />
<br />
Reviewing the examples from the above figure, it can be observed that the method from this paper repeats fewer local features from the image than a competing feed forward network method [8]. While the analysis is qualitative, the authors claim that their method produces "more visually pleasing results".<br />
<br />
=Conclusion=<br />
Only a couple years ago were CNNs first used to stylize images. Today, a host of improvements have been developed, optimizing the original work of Gatys et al. for a number of different situations. Using additional training per style image, computational efficiency and image quality can be increased. However, the trained network then depends on that specific style image, or in some cases such as in [3], a set of style images. Till now, limited work has taken place in improving Gatys et al.'s method for arbitrary style images. The authors of this paper developed and evaluated a novel method for arbitrary style transfer. Their method and Gatys et al.'s method share the use of a VGG-19 CNN as the initial processing step. However, the authors replaced iterative optimization with whitening and colour transforms, which can be applied in a single step. This yields a decrease in computational time while maintaining generality with respect to the style image. After their CNN auto-encoder is initially trained no further training is required. This allows their method to be style agnostic. Their method also performs favourably, in terms of image quality, when compared to other current work.<br />
<br />
=Critique=<br />
In the paper, the authors only experimented with layers of VGG19. Given that architectures such as ResNet and Xception perform better on image recognition tasks, it would be interesting to see how residual layers and/or Inception modules may be applied to the task of disentangling style and content and whether they would improve performance relative to the results presented in the current paper is the encoder used were to utilize layers from these alternative convolutional architectures. Additionally, it is worth exploring whether one can invent a probabilistic and/or generative version of the encoder-decoder architecture used in the paper. More precisely, is it possible to come up with something in the spirit of variational autoencoders, wherein we the bottleneck layer can be used to sample noise vectors, which can then be input into each of the decoder units to generate synthetic style and content images.<br />
Alternative attempts would also involve the study of generative adversarial networks with a perturbation threshold value. GANs can produce surreal images, where the underlying structure (content) is preserved ( in CNNs the filters learn the edges and surfaces and shape of the image), provided the Discriminator is trained for style classification ( training set consists of images pertaining the style that requires to be transferred).<br />
<br />
=References=<br />
[1] L. A. Gatys, A. S. Ecker, and M. Bethge. Texture synthesis using convolutional neural networks. In NIPS, 2015.<br />
<br />
[2] L. A. Gatys, A. S. Ecker, and M. Bethge. Image style transfer using convolutional neural networks. In CVPR, 2016.<br />
<br />
[3] V. Dumoulin, J. Shlens, and M. Kudlur. A learned representation for artistic style. In ICLR, 2017.<br />
<br />
[4] J. Johnson, A. Alahi, and L. Fei-Fei. Perceptual losses for real-time style transfer and super-resolution. In ECCV, 2016<br />
<br />
[5] R. Picard. MAS 622J/1.126J: Pattern Recognition and Analysis, Lecture 4. http://courses.media.mit.edu/2010fall/mas622j/whiten.pdf<br />
<br />
[6] T. Q. Chen and M. Schmidt. Fast patch-based style transfer of arbitrary style. arXiv preprint arXiv:1612.04337, 2016.<br />
<br />
[7] X. Huang and S. Belongie. Arbitrary style transfer in real-time with adaptive instance normalization. arXiv preprint arXiv:1703.06868, 2017.<br />
<br />
[8] D. Ulyanov, V. Lebedev, A. Vedaldi, and V. Lempitsky. Texture networks: Feed-forward synthesis of textures and stylized images. In ICML, 2016.<br />
<br />
[9] Leon A. Gatys, Alexander S. Ecker, Matthias Bethge, A Neural Algorithm of Artistic Style, https://arxiv.org/abs/1508.06576<br />
<br />
[10] Karen Simonyan et al. Very Deep Convolutional Networks for Large-Scale Image Recognition<br />
<br />
[11] VGG Architectures - [http://www.robots.ox.ac.uk/~vgg/research/very_deep/| More Details]<br />
<br />
[12] Mechrez, R., Shechtman, E., & Zelnik-Manor, L. (2017). Photorealistic Style Transfer with Screened Poisson Equation. arXiv preprint arXiv:1709.09828.<br />
<br />
[13] Y. Li, C. Fang, J. Yang, Z. Wang, X. Lu, and M.-H. Yang. Diversified texture synthesis with feed-forward networks. In CVPR, 2017<br />
<br />
[14] D. Chen, L. Yuan, J. Liao, N. Yu, and G. Hua. Stylebank: An explicit representation for neural image style transfer. In CVPR, 2017<br />
<br />
Implementation Example: https://github.com/titu1994/Neural-Style-Transfer</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Universal_Style_Transfer_via_Feature_Transforms&diff=30828Universal Style Transfer via Feature Transforms2017-11-19T19:14:12Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>=Introduction=<br />
When viewing an image, whether it is a photograph or a painting, two types of mutually exclusive data are present. First, there is the content of the image, such as a person in a portrait. However, the content does not uniquely define the image. Consider a case where multiple artists paint a portrait of an identical subject, the results would vary despite the content being invariant. The cause of the variance is rooted in the style of each particular artist. Therefore, style transfer between two images results in the content being unaffected but the style being copied. Style transfer is an important image editing task which enables the creation of new artistic works. Typically one image is termed the content/reference image, whose style is discarded. The other image is called the style image, whose style, but the not content is copied to the content image.<br />
<br />
Deep learning techniques have been shown to be effective methods for implementing style transfer. Previous methods have been successful but with several key limitations and often trade off between generalization, quality and efficiency. Either they are fast, but have very few styles that can be transferred or they can handle arbitrary styles but are no longer efficient. The presented paper establishes a compromise between these two extremes by using only whitening and coloring transforms (WCT) to transfer a style within a feedforward image reconstruction architecture. No training of the underlying deep network is required per style.<br />
<br />
=Related Work=<br />
Gatys et al. developed a new method for generating textures from sample images in 2015 [1] and extended their approach to style transfer by 2016 [2]. They proposed the use of a pre-trained convolutional neural network (CNN) to separate content and style of input images. Having proven successful, a number of improvements quickly developed, reducing computational time, increasing the diversity of transferrable styes, and improving the quality of the results. Central to these approaches and of the present paper is the use of a CNN.<br />
<br />
In 2017, Mechrez et al. [12] proposed an approach that takes as input a stylized image and makes it more photorealistic. Their approach relied on the Screened Poisson Equation, maintaining the fidelity of the stylized image while constraining the gradients to those of the original input image. The method they proposed was fast, simple, fully automatic and showed positive progress in making a stylized image photorealistic.<br />
<br />
Alternative attempts, by using a single network to transfer<br />
multiple styles, include models conditioned on binary selection units [13], a network that learns a set of new filters for every new style [15], and a novel conditional normalization layer that learns normalization parameters for each style [3]<br />
==How Content and Style are Extracted using CNNs==<br />
A CNN was chosen due to its ability to extract high level feature from images. These features can be interpreted in two ways. Within layer <math> l </math> there are <math> N_l </math> feature maps of size <math> M_l </math>. With a particular input image, the feature maps are given by <math> F_{i,j}^l </math> where <math> i </math> and <math> j </math> locate the map within the layer. Starting with a white noise image and an reference (content) image, the features can be transferred by minimizing<br />
<br />
<center><br />
<math> \mathcal{L}_{content} = \frac{1}{2} \sum_{i,j} \left( F_{i,j}^l - P_{i,j}^l \right)^2 </math><br />
</center><br />
<br />
where <math> P_{i,j} </math> denotes the feature map output caused by the white noise image. Therefore this loss function preserves the content of the reference image. The style is described using a Gram matrix given by<br />
<br />
<center><br />
<math><br />
G_{i,j}^l = \sum_k F_{i,k}^l F_{j,k}^l<br />
</math><br />
</center><br />
<br />
Gram matrix $G$ of a set of vectors $v_1,\dots,v_n$ is the matrix of all possible inner products whose entries are given by $G_{ij}=v_i^Tv_j$. The loss function that describes a difference in style between two images is equal to:<br />
<br />
<center><br />
<math><br />
\mathcal{L}_{style} = \frac{1}{4 N_l^2 M_l^2} \sum_{i,j} \left(G_{i,j}^l - A_{i,j}^l \right)^2<br />
</math><br />
</center><br />
<br />
where <math> A_{i,j}^l </math> and <math> G_{i,j}^l </math> are the Gram matrices of the generated image and style image respectively. Therefore three images are required, a style image, a content image and an initial white noise image. Iterative optimization is then used to add content from one image to the white noise image, and style from the other. An additional parameter is used to balance the ratio of these loss functions.<br />
<br />
The 19-layer ImageNet trained VGG network was chosen by Gatys et al. VGG-19 is still commonly used in more recent works as will be shown in the presented paper, although training datasets vary. Such CNNs are typically used in classification problems by finalizing their output through a series of full connected layers. For content and style extraction it is the convolutional layers that are required. The method of Gatys et al. is style independent, since the CNN does not need to be trained for each style image. However the process of iterative optimization to generate the output image is computational expensive.<br />
<br />
==Other Methods==<br />
Other methods avoid the inefficiency of iterative optimization by training a network/networks on a set of styles. The network then directly transfers the style from the style image to the content image without solving the iterative optimization problem. V. Dumoulin et al. trained a single network on $N$ styles [3]. This improved upon previous work where a network was required per style [4]. The stylized output image was generated by simply running a feedforward pass of the network on the content image. While efficiency is high, the method is no longer able to apply an arbitrary style without retraining.<br />
<br />
=Methodology=<br />
Li et al. have proposed a novel method for generating the stylized image. A CNN is still used as in Gatys et al. to extract content and style. However, the stylized image is not generated through iterative optimization or a feed-forward pass as required by previous methods. Instead, whitening and colour transforms are used.<br />
<br />
==Image Reconstruction==<br />
[[File:image_resconstruction.png|thumb|150px|right|alt=Training a single decoder.|Training a single decoder. X denotes the layer of the VGG encoder that the decoder receives as input.]]<br />
An auto-encoder network is used to first encode an input image into a set of feature maps, and then decode it back to an image as shown in the adjacent figure. The encoder network used is VGG-19. This network is reponsible for obtaining feature maps (similar to Gatys et al.). The output of each of the first five layers is then fed into a corresponding decoder network, which is a mirrored version of VGG-19. Each decoder network then decodes the feature maps of the $l$th layer producing an output image. A mechanism for transferring style will be implemented by manipulating the feature maps between the encoder and decoder networks.<br />
<br />
First, the auto-encoder network needs to be trained. The following loss function is used<br />
<br />
<center><br />
<math><br />
<br />
\mathcal{L} = || I_{output} - I_{input} ||_2^2 + \lambda || \Phi(I_{output}) - \Phi(I_{input})||_2^2<br />
<br />
</math><br />
</center><br />
<br />
where $I_{input}$ and $I_{output}$ are the input and output images of the auto-encoder. $\Phi$ is the VGG encoder. The first term of the loss is the pixel reconstruction loss, while the second term is feature loss. Recall from "Related Work" that the feature maps correspond to the content of the image. Therefore the second term can also be seen as penalising for content differences that arise due to the encoder network. The network was trained using the Microsoft COCO dataset. <br />
<br />
They use whitening and coloring transforms to directly transform the $f_c$ (VGG feature<br />
map of content image at a certain layer) to match the covariance matrix of $f_s$ (VGG feature<br />
map of style image). This process is consisted of two steps, i.e., whitening and coloring transform. Note that the decoder will reconstruct the original content image if $f_c$ is directly fed into it.<br />
<br />
==Whitening Transform==<br />
Whitening first requires that the covariance of the data is a diagonal matrix. This is done by solving for the covariance matrix's eigenvalues and eigenvector matrices. Whitening then forces the diagonal elements of the eigenvalue matrix to be the same. This is achieved for a feature map from VGG through the following steps.<br />
<br />
# The feature map $f_c$ is extracted from a layer of the encoder network after activation on the content image. This is the data to be whitened.<br />
# $f_c$ is centered by subtracting its mean vector $m_c$.<br />
# Then, the eigenvectors $E_c$ and eigenvalues $D_c$ are found for the covariance matrix of $f_c$. <br />
# The whitened feature map is then given by $\hat{f}_c = E_c D_c^{-1/2} E_c^T f_c$. <br />
<br />
If interested, the derivation of the whitening equation can be seen in [5]. Li et al. found that whitening removed styles from the image.<br />
<br />
==Colour Transform==<br />
However, whitening does not transfer style from the style image. It only uses feature maps from the content image. The colour transform uses both $\hat{f}_c$ from above and $f_s$, the feature map from the style image.<br />
<br />
# $f_s$ is centered by subtracting its mean vector $m_s$.<br />
# Then, the eigenvectors $E_s$ and eigenvalues $D_s$ are calculated for the covariance matrix of $f_s$.<br />
# The colour transform is given by $\hat{f}_{cs} = E_s D_s^{1/2} E_s^T \hat{f}_c$.<br />
# Recenter $\hat{f}_{cs}$ using $m_s$.<br />
<br />
Intuitively, colouring results in a correlation between the $\hat{f}_c$ and $f_s$ feature maps. This is where the style transfer takes place.<br />
<br />
==Content/Style Balance==<br />
Using just $\hat{f}_{cs}$ as the input to the decoder may create a result that is too extreme in style. To balance content and style a new parameter $\alpha$ is defined.<br />
<br />
<center><br />
<math><br />
<br />
\hat{f}_{cs} = \alpha \hat{f}_{cs} + (1 - \alpha) f_c<br />
<br />
</math><br />
</center><br />
<br />
==Using Multiple Layers==<br />
It has been previously mentioned that multiple decoders were trained, one for each of the first five layers of the encoder network. Each layer of a CNN perceives features at different levels. Levels close to the input image will detect lower level local features such as edges. Those levels deeper into the network will detect more complex global features. The style transfer algorithm is applied at each of these levels, which yields the question as to which results, as shown below, to use.<br />
<br />
[[File:multilevel_features.png|thumb|700px|center|alt=Results of style transfer from each of the first five layers of the encoder network.|Results of style transfer from each of the first five layers of the encoder network.]]<br />
<br />
Ideally, the results of each layer should be used to build the final output image. This captures the entire range of features detected by the encoder network. First, one full pass of the network is performed. Then the stylised image from the deepest layer (Relu_5_1 in this case) is taken and used as the content image for another iteration of the algorithm, where then the next layer (Relu_4_1) is used as the output. These steps are repeated until the final image is produced from the shallowest layer. This process is summarised in the figure below.<br />
<br />
[[File:process_summary.png|thumb|700px|center|alt=Process summary of the multi-level stylization algorithm.|The content (C) and style (S) are fed to the VGG encoding network. The output image (I) after a whitening and colour transform (WCT) is taken from the deepest level's decoder. The process is iteratively repeated until the most shallow layer is reached.]]<br />
<br />
The authors note that the transformations must be applied first at the highest level (most abstract) layers, which capture complicated local structures and pass this transformed image to lower layers, which improve on details. They observe that reversing this order (lowest to highest) leads to images with low visual quality, as low-level information cannot be preserved after manipulating high level features.<br />
<br />
[[File:Universal_Style_Transfer_Coarse_to_Fine.JPG|thumb|700px|center|alt=(a)-(c) Output from intermediate layers. (d) Reversed transformation order.|(a)-(c) Output from intermediate layers. (d) Reversed transformation order.]]<br />
<br />
=Evaluation=<br />
The success of style transfer might appear hard to quantify as it relies on qualitative judgement. However, the extremes of transferring no style, or transferring only style can be considered as performing poorly. Consistent transfer of style throughout the entire image is another parameter of success. Ideally, the viewer can recognize the content of the image, while seeing it expressed in an alternative style. Quantitatively, the quality of the style transfer can be calculated by taking the covariance matrix difference $L_s$ between the resulting image and the original style. The results of the presented paper also need to be considered within the contexts of generality, efficiency and training requirements.<br />
<br />
==Style Transfer==<br />
A number of style transfer examples are presented relative to other works. <br />
<br />
[[File:transfer_results_label.jpg|thumb|700px|center|alt=Style transfer results of the presented paper.|A: See [6]. B: See [7]. C: See [8]. D: Gatys et al. iterative optimization, see [2]. E: This paper's results.]]<br />
<br />
Li et al. then obtained the average $L_s$ using 10 random content images across 40 style images. They had the lowest average $log(L_s)$ of all referenced works at 6.3. Next lowest was Gatys et al. [2] with $log(L_s) = 6.7$. It should be noted that while $L_s$ quantitatively calculates the success of the style transfer, results are still subject to the viewer's impression. Reviewing the transfer results, rows five and six for Gatys et al.'s method shows local minimization issues. However, their method still achieves a competitive $L_s$ score.<br />
<br />
==Transfer Efficiency==<br />
It was hypothesized by Li et al. that using WCT would enable faster run-times than [2] while still supporting arbitrary style transfer. For a 256x256 image, using a 12GB TITAN X, they achieved a transfer time of 1.5 seconds. Gatys et al.'s method [2] required 21.2 seconds. The pure feed-forward approaches [7], and [8] had times equal to or less than 0.2 seconds. [6] had a time comparable to the presented paper's method. However, [6,7,8] do not generalize well to multiple styles as training is required. Therefore this paper obtained a near 15x speed up for a style agnostic transfer algorithm when compared to leading previous work. The authors also note that WCT was done using the CPU. They intend to port WCT to the GPU and expect to see the computational time be further reduced.<br />
<br />
==Other Applications==<br />
Li et al.'s method can also be used for texture synthesis. This was the original work of Gatys et. al. before they applied their algorithm to style transfer problems. Texture synthesis takes a reference texture/image and creates new textures from it. With proper boundary conditions enforced these synthesized textures can be tileable. Alternatively, higher resolution textures can be generated. Texture synthesis has applications in areas such as computer graphics, allowing for large surfaces to be texture mapped.<br />
<br />
The content image is set as white noise, similar to how [2] initializes their output image. Then the reference texture/image is set as the style image. Since the content image is initially random white noise, then the features generated by the encoder of this image are also random. Li et al. state that this increases the diversity of the resulting output textures.<br />
<br />
[[File:texture_synthesis_label.jpg|thumb|700px|center|alt=Texture synthesis results.|A: Reference image/texture. B: Result from [8]. C: Result of present paper.]]<br />
<br />
Reviewing the examples from the above figure, it can be observed that the method from this paper repeats fewer local features from the image than a competing feed forward network method [8]. While the analysis is qualitative, the authors claim that their method produces "more visually pleasing results".<br />
<br />
=Conclusion=<br />
Only a couple years ago were CNNs first used to stylize images. Today, a host of improvements have been developed, optimizing the original work of Gatys et al. for a number of different situations. Using additional training per style image, computational efficiency and image quality can be increased. However, the trained network then depends on that specific style image, or in some cases such as in [3], a set of style images. Till now, limited work has taken place in improving Gatys et al.'s method for arbitrary style images. The authors of this paper developed and evaluated a novel method for arbitrary style transfer. Their method and Gatys et al.'s method share the use of a VGG-19 CNN as the initial processing step. However, the authors replaced iterative optimization with whitening and colour transforms, which can be applied in a single step. This yields a decrease in computational time while maintaining generality with respect to the style image. After their CNN auto-encoder is initially trained no further training is required. This allows their method to be style agnostic. Their method also performs favourably, in terms of image quality, when compared to other current work.<br />
<br />
=Critique=<br />
In the paper, the authors only experimented with layers of VGG19. Given that architectures such as ResNet and Xception perform better on image recognition tasks, it would be interesting to see how residual layers and/or Inception modules may be applied to the task of disentangling style and content and whether they would improve performance relative to the results presented in the current paper is the encoder used were to utilize layers from these alternative convolutional architectures. Additionally, it is worth exploring whether one can invent a probabilistic and/or generative version of the encoder-decoder architecture used in the paper. More precisely, is it possible to come up with something in the spirit of variational autoencoders, wherein we the bottleneck layer can be used to sample noise vectors, which can then be input into each of the decoder units to generate synthetic style and content images.<br />
Alternative attempts would also involve the study of generative adversarial networks with a perturbation threshold value. GANs can produce surreal images, where the underlying structure (content) is preserved ( in CNNs the filters learn the edges and surfaces and shape of the image), provided the Discriminator is trained for style classification ( training set consists of images pertaining the style that requires to be transferred).<br />
<br />
=References=<br />
[1] L. A. Gatys, A. S. Ecker, and M. Bethge. Texture synthesis using convolutional neural networks. In NIPS, 2015.<br />
<br />
[2] L. A. Gatys, A. S. Ecker, and M. Bethge. Image style transfer using convolutional neural networks. In CVPR, 2016.<br />
<br />
[3] V. Dumoulin, J. Shlens, and M. Kudlur. A learned representation for artistic style. In ICLR, 2017.<br />
<br />
[4] J. Johnson, A. Alahi, and L. Fei-Fei. Perceptual losses for real-time style transfer and super-resolution. In ECCV, 2016<br />
<br />
[5] R. Picard. MAS 622J/1.126J: Pattern Recognition and Analysis, Lecture 4. http://courses.media.mit.edu/2010fall/mas622j/whiten.pdf<br />
<br />
[6] T. Q. Chen and M. Schmidt. Fast patch-based style transfer of arbitrary style. arXiv preprint arXiv:1612.04337, 2016.<br />
<br />
[7] X. Huang and S. Belongie. Arbitrary style transfer in real-time with adaptive instance normalization. arXiv preprint arXiv:1703.06868, 2017.<br />
<br />
[8] D. Ulyanov, V. Lebedev, A. Vedaldi, and V. Lempitsky. Texture networks: Feed-forward synthesis of textures and stylized images. In ICML, 2016.<br />
<br />
[9] Leon A. Gatys, Alexander S. Ecker, Matthias Bethge, A Neural Algorithm of Artistic Style, https://arxiv.org/abs/1508.06576<br />
<br />
[10] Karen Simonyan et al. Very Deep Convolutional Networks for Large-Scale Image Recognition<br />
<br />
[11] VGG Architectures - [http://www.robots.ox.ac.uk/~vgg/research/very_deep/| More Details]<br />
<br />
[12] Mechrez, R., Shechtman, E., & Zelnik-Manor, L. (2017). Photorealistic Style Transfer with Screened Poisson Equation. arXiv preprint arXiv:1709.09828.<br />
<br />
[13] Y. Li, C. Fang, J. Yang, Z. Wang, X. Lu, and M.-H. Yang. Diversified texture synthesis with feed-forward networks. In CVPR, 2017<br />
<br />
[14] D. Chen, L. Yuan, J. Liao, N. Yu, and G. Hua. Stylebank: An explicit representation for neural image style transfer. In CVPR, 2017<br />
<br />
Implementation Example: https://github.com/titu1994/Neural-Style-Transfer</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Universal_Style_Transfer_via_Feature_Transforms&diff=30825Universal Style Transfer via Feature Transforms2017-11-19T19:12:44Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>=Introduction=<br />
When viewing an image, whether it is a photograph or a painting, two types of mutually exclusive data are present. First, there is the content of the image, such as a person in a portrait. However, the content does not uniquely define the image. Consider a case where multiple artists paint a portrait of an identical subject, the results would vary despite the content being invariant. The cause of the variance is rooted in the style of each particular artist. Therefore, style transfer between two images results in the content being unaffected but the style being copied. Style transfer is an important image editing task which enables the creation of new artistic works. Typically one image is termed the content/reference image, whose style is discarded. The other image is called the style image, whose style, but the not content is copied to the content image.<br />
<br />
Deep learning techniques have been shown to be effective methods for implementing style transfer. Previous methods have been successful but with several key limitations. Either they are fast, but have very few styles that can be transferred or they can handle arbitrary styles but are no longer efficient. The presented paper establishes a compromise between these two extremes by using only whitening and coloring transforms (WCT) to transfer a style within a feedforward image reconstruction architecture. No training of the underlying deep network is required per style.<br />
<br />
=Related Work=<br />
Gatys et al. developed a new method for generating textures from sample images in 2015 [1] and extended their approach to style transfer by 2016 [2]. They proposed the use of a pre-trained convolutional neural network (CNN) to separate content and style of input images. Having proven successful, a number of improvements quickly developed, reducing computational time, increasing the diversity of transferrable styes, and improving the quality of the results. Central to these approaches and of the present paper is the use of a CNN.<br />
<br />
In 2017, Mechrez et al. [12] proposed an approach that takes as input a stylized image and makes it more photorealistic. Their approach relied on the Screened Poisson Equation, maintaining the fidelity of the stylized image while constraining the gradients to those of the original input image. The method they proposed was fast, simple, fully automatic and showed positive progress in making a stylized image photorealistic.<br />
<br />
Alternative attempts, by using a single network to transfer<br />
multiple styles, include models conditioned on binary selection units [13], a network that learns a set of new filters for every new style [15], and a novel conditional normalization layer that learns normalization parameters for each style [3]<br />
==How Content and Style are Extracted using CNNs==<br />
A CNN was chosen due to its ability to extract high level feature from images. These features can be interpreted in two ways. Within layer <math> l </math> there are <math> N_l </math> feature maps of size <math> M_l </math>. With a particular input image, the feature maps are given by <math> F_{i,j}^l </math> where <math> i </math> and <math> j </math> locate the map within the layer. Starting with a white noise image and an reference (content) image, the features can be transferred by minimizing<br />
<br />
<center><br />
<math> \mathcal{L}_{content} = \frac{1}{2} \sum_{i,j} \left( F_{i,j}^l - P_{i,j}^l \right)^2 </math><br />
</center><br />
<br />
where <math> P_{i,j} </math> denotes the feature map output caused by the white noise image. Therefore this loss function preserves the content of the reference image. The style is described using a Gram matrix given by<br />
<br />
<center><br />
<math><br />
G_{i,j}^l = \sum_k F_{i,k}^l F_{j,k}^l<br />
</math><br />
</center><br />
<br />
Gram matrix $G$ of a set of vectors $v_1,\dots,v_n$ is the matrix of all possible inner products whose entries are given by $G_{ij}=v_i^Tv_j$. The loss function that describes a difference in style between two images is equal to:<br />
<br />
<center><br />
<math><br />
\mathcal{L}_{style} = \frac{1}{4 N_l^2 M_l^2} \sum_{i,j} \left(G_{i,j}^l - A_{i,j}^l \right)^2<br />
</math><br />
</center><br />
<br />
where <math> A_{i,j}^l </math> and <math> G_{i,j}^l </math> are the Gram matrices of the generated image and style image respectively. Therefore three images are required, a style image, a content image and an initial white noise image. Iterative optimization is then used to add content from one image to the white noise image, and style from the other. An additional parameter is used to balance the ratio of these loss functions.<br />
<br />
The 19-layer ImageNet trained VGG network was chosen by Gatys et al. VGG-19 is still commonly used in more recent works as will be shown in the presented paper, although training datasets vary. Such CNNs are typically used in classification problems by finalizing their output through a series of full connected layers. For content and style extraction it is the convolutional layers that are required. The method of Gatys et al. is style independent, since the CNN does not need to be trained for each style image. However the process of iterative optimization to generate the output image is computational expensive.<br />
<br />
==Other Methods==<br />
Other methods avoid the inefficiency of iterative optimization by training a network/networks on a set of styles. The network then directly transfers the style from the style image to the content image without solving the iterative optimization problem. V. Dumoulin et al. trained a single network on $N$ styles [3]. This improved upon previous work where a network was required per style [4]. The stylized output image was generated by simply running a feedforward pass of the network on the content image. While efficiency is high, the method is no longer able to apply an arbitrary style without retraining.<br />
<br />
=Methodology=<br />
Li et al. have proposed a novel method for generating the stylized image. A CNN is still used as in Gatys et al. to extract content and style. However, the stylized image is not generated through iterative optimization or a feed-forward pass as required by previous methods. Instead, whitening and colour transforms are used.<br />
<br />
==Image Reconstruction==<br />
[[File:image_resconstruction.png|thumb|150px|right|alt=Training a single decoder.|Training a single decoder. X denotes the layer of the VGG encoder that the decoder receives as input.]]<br />
An auto-encoder network is used to first encode an input image into a set of feature maps, and then decode it back to an image as shown in the adjacent figure. The encoder network used is VGG-19. This network is reponsible for obtaining feature maps (similar to Gatys et al.). The output of each of the first five layers is then fed into a corresponding decoder network, which is a mirrored version of VGG-19. Each decoder network then decodes the feature maps of the $l$th layer producing an output image. A mechanism for transferring style will be implemented by manipulating the feature maps between the encoder and decoder networks.<br />
<br />
First, the auto-encoder network needs to be trained. The following loss function is used<br />
<br />
<center><br />
<math><br />
<br />
\mathcal{L} = || I_{output} - I_{input} ||_2^2 + \lambda || \Phi(I_{output}) - \Phi(I_{input})||_2^2<br />
<br />
</math><br />
</center><br />
<br />
where $I_{input}$ and $I_{output}$ are the input and output images of the auto-encoder. $\Phi$ is the VGG encoder. The first term of the loss is the pixel reconstruction loss, while the second term is feature loss. Recall from "Related Work" that the feature maps correspond to the content of the image. Therefore the second term can also be seen as penalising for content differences that arise due to the encoder network. The network was trained using the Microsoft COCO dataset. <br />
<br />
They use whitening and coloring transforms to directly transform the $f_c$ (VGG feature<br />
map of content image at a certain layer) to match the covariance matrix of $f_s$ (VGG feature<br />
map of style image). This process is consisted of two steps, i.e., whitening and coloring transform. Note that the decoder will reconstruct the original content image if $f_c$ is directly fed into it.<br />
<br />
==Whitening Transform==<br />
Whitening first requires that the covariance of the data is a diagonal matrix. This is done by solving for the covariance matrix's eigenvalues and eigenvector matrices. Whitening then forces the diagonal elements of the eigenvalue matrix to be the same. This is achieved for a feature map from VGG through the following steps.<br />
<br />
# The feature map $f_c$ is extracted from a layer of the encoder network after activation on the content image. This is the data to be whitened.<br />
# $f_c$ is centered by subtracting its mean vector $m_c$.<br />
# Then, the eigenvectors $E_c$ and eigenvalues $D_c$ are found for the covariance matrix of $f_c$. <br />
# The whitened feature map is then given by $\hat{f}_c = E_c D_c^{-1/2} E_c^T f_c$. <br />
<br />
If interested, the derivation of the whitening equation can be seen in [5]. Li et al. found that whitening removed styles from the image.<br />
<br />
==Colour Transform==<br />
However, whitening does not transfer style from the style image. It only uses feature maps from the content image. The colour transform uses both $\hat{f}_c$ from above and $f_s$, the feature map from the style image.<br />
<br />
# $f_s$ is centered by subtracting its mean vector $m_s$.<br />
# Then, the eigenvectors $E_s$ and eigenvalues $D_s$ are calculated for the covariance matrix of $f_s$.<br />
# The colour transform is given by $\hat{f}_{cs} = E_s D_s^{1/2} E_s^T \hat{f}_c$.<br />
# Recenter $\hat{f}_{cs}$ using $m_s$.<br />
<br />
Intuitively, colouring results in a correlation between the $\hat{f}_c$ and $f_s$ feature maps. This is where the style transfer takes place.<br />
<br />
==Content/Style Balance==<br />
Using just $\hat{f}_{cs}$ as the input to the decoder may create a result that is too extreme in style. To balance content and style a new parameter $\alpha$ is defined.<br />
<br />
<center><br />
<math><br />
<br />
\hat{f}_{cs} = \alpha \hat{f}_{cs} + (1 - \alpha) f_c<br />
<br />
</math><br />
</center><br />
<br />
==Using Multiple Layers==<br />
It has been previously mentioned that multiple decoders were trained, one for each of the first five layers of the encoder network. Each layer of a CNN perceives features at different levels. Levels close to the input image will detect lower level local features such as edges. Those levels deeper into the network will detect more complex global features. The style transfer algorithm is applied at each of these levels, which yields the question as to which results, as shown below, to use.<br />
<br />
[[File:multilevel_features.png|thumb|700px|center|alt=Results of style transfer from each of the first five layers of the encoder network.|Results of style transfer from each of the first five layers of the encoder network.]]<br />
<br />
Ideally, the results of each layer should be used to build the final output image. This captures the entire range of features detected by the encoder network. First, one full pass of the network is performed. Then the stylised image from the deepest layer (Relu_5_1 in this case) is taken and used as the content image for another iteration of the algorithm, where then the next layer (Relu_4_1) is used as the output. These steps are repeated until the final image is produced from the shallowest layer. This process is summarised in the figure below.<br />
<br />
[[File:process_summary.png|thumb|700px|center|alt=Process summary of the multi-level stylization algorithm.|The content (C) and style (S) are fed to the VGG encoding network. The output image (I) after a whitening and colour transform (WCT) is taken from the deepest level's decoder. The process is iteratively repeated until the most shallow layer is reached.]]<br />
<br />
The authors note that the transformations must be applied first at the highest level (most abstract) layers, which capture complicated local structures and pass this transformed image to lower layers, which improve on details. They observe that reversing this order (lowest to highest) leads to images with low visual quality, as low-level information cannot be preserved after manipulating high level features.<br />
<br />
[[File:Universal_Style_Transfer_Coarse_to_Fine.JPG|thumb|700px|center|alt=(a)-(c) Output from intermediate layers. (d) Reversed transformation order.|(a)-(c) Output from intermediate layers. (d) Reversed transformation order.]]<br />
<br />
=Evaluation=<br />
The success of style transfer might appear hard to quantify as it relies on qualitative judgement. However, the extremes of transferring no style, or transferring only style can be considered as performing poorly. Consistent transfer of style throughout the entire image is another parameter of success. Ideally, the viewer can recognize the content of the image, while seeing it expressed in an alternative style. Quantitatively, the quality of the style transfer can be calculated by taking the covariance matrix difference $L_s$ between the resulting image and the original style. The results of the presented paper also need to be considered within the contexts of generality, efficiency and training requirements.<br />
<br />
==Style Transfer==<br />
A number of style transfer examples are presented relative to other works. <br />
<br />
[[File:transfer_results_label.jpg|thumb|700px|center|alt=Style transfer results of the presented paper.|A: See [6]. B: See [7]. C: See [8]. D: Gatys et al. iterative optimization, see [2]. E: This paper's results.]]<br />
<br />
Li et al. then obtained the average $L_s$ using 10 random content images across 40 style images. They had the lowest average $log(L_s)$ of all referenced works at 6.3. Next lowest was Gatys et al. [2] with $log(L_s) = 6.7$. It should be noted that while $L_s$ quantitatively calculates the success of the style transfer, results are still subject to the viewer's impression. Reviewing the transfer results, rows five and six for Gatys et al.'s method shows local minimization issues. However, their method still achieves a competitive $L_s$ score.<br />
<br />
==Transfer Efficiency==<br />
It was hypothesized by Li et al. that using WCT would enable faster run-times than [2] while still supporting arbitrary style transfer. For a 256x256 image, using a 12GB TITAN X, they achieved a transfer time of 1.5 seconds. Gatys et al.'s method [2] required 21.2 seconds. The pure feed-forward approaches [7], and [8] had times equal to or less than 0.2 seconds. [6] had a time comparable to the presented paper's method. However, [6,7,8] do not generalize well to multiple styles as training is required. Therefore this paper obtained a near 15x speed up for a style agnostic transfer algorithm when compared to leading previous work. The authors also note that WCT was done using the CPU. They intend to port WCT to the GPU and expect to see the computational time be further reduced.<br />
<br />
==Other Applications==<br />
Li et al.'s method can also be used for texture synthesis. This was the original work of Gatys et. al. before they applied their algorithm to style transfer problems. Texture synthesis takes a reference texture/image and creates new textures from it. With proper boundary conditions enforced these synthesized textures can be tileable. Alternatively, higher resolution textures can be generated. Texture synthesis has applications in areas such as computer graphics, allowing for large surfaces to be texture mapped.<br />
<br />
The content image is set as white noise, similar to how [2] initializes their output image. Then the reference texture/image is set as the style image. Since the content image is initially random white noise, then the features generated by the encoder of this image are also random. Li et al. state that this increases the diversity of the resulting output textures.<br />
<br />
[[File:texture_synthesis_label.jpg|thumb|700px|center|alt=Texture synthesis results.|A: Reference image/texture. B: Result from [8]. C: Result of present paper.]]<br />
<br />
Reviewing the examples from the above figure, it can be observed that the method from this paper repeats fewer local features from the image than a competing feed forward network method [8]. While the analysis is qualitative, the authors claim that their method produces "more visually pleasing results".<br />
<br />
=Conclusion=<br />
Only a couple years ago were CNNs first used to stylize images. Today, a host of improvements have been developed, optimizing the original work of Gatys et al. for a number of different situations. Using additional training per style image, computational efficiency and image quality can be increased. However, the trained network then depends on that specific style image, or in some cases such as in [3], a set of style images. Till now, limited work has taken place in improving Gatys et al.'s method for arbitrary style images. The authors of this paper developed and evaluated a novel method for arbitrary style transfer. Their method and Gatys et al.'s method share the use of a VGG-19 CNN as the initial processing step. However, the authors replaced iterative optimization with whitening and colour transforms, which can be applied in a single step. This yields a decrease in computational time while maintaining generality with respect to the style image. After their CNN auto-encoder is initially trained no further training is required. This allows their method to be style agnostic. Their method also performs favourably, in terms of image quality, when compared to other current work.<br />
<br />
=Critique=<br />
In the paper, the authors only experimented with layers of VGG19. Given that architectures such as ResNet and Xception perform better on image recognition tasks, it would be interesting to see how residual layers and/or Inception modules may be applied to the task of disentangling style and content and whether they would improve performance relative to the results presented in the current paper is the encoder used were to utilize layers from these alternative convolutional architectures. Additionally, it is worth exploring whether one can invent a probabilistic and/or generative version of the encoder-decoder architecture used in the paper. More precisely, is it possible to come up with something in the spirit of variational autoencoders, wherein we the bottleneck layer can be used to sample noise vectors, which can then be input into each of the decoder units to generate synthetic style and content images.<br />
Alternative attempts would also involve the study of generative adversarial networks with a perturbation threshold value. GANs can produce surreal images, where the underlying structure (content) is preserved ( in CNNs the filters learn the edges and surfaces and shape of the image), provided the Discriminator is trained for style classification ( training set consists of images pertaining the style that requires to be transferred).<br />
<br />
=References=<br />
[1] L. A. Gatys, A. S. Ecker, and M. Bethge. Texture synthesis using convolutional neural networks. In NIPS, 2015.<br />
<br />
[2] L. A. Gatys, A. S. Ecker, and M. Bethge. Image style transfer using convolutional neural networks. In CVPR, 2016.<br />
<br />
[3] V. Dumoulin, J. Shlens, and M. Kudlur. A learned representation for artistic style. In ICLR, 2017.<br />
<br />
[4] J. Johnson, A. Alahi, and L. Fei-Fei. Perceptual losses for real-time style transfer and super-resolution. In ECCV, 2016<br />
<br />
[5] R. Picard. MAS 622J/1.126J: Pattern Recognition and Analysis, Lecture 4. http://courses.media.mit.edu/2010fall/mas622j/whiten.pdf<br />
<br />
[6] T. Q. Chen and M. Schmidt. Fast patch-based style transfer of arbitrary style. arXiv preprint arXiv:1612.04337, 2016.<br />
<br />
[7] X. Huang and S. Belongie. Arbitrary style transfer in real-time with adaptive instance normalization. arXiv preprint arXiv:1703.06868, 2017.<br />
<br />
[8] D. Ulyanov, V. Lebedev, A. Vedaldi, and V. Lempitsky. Texture networks: Feed-forward synthesis of textures and stylized images. In ICML, 2016.<br />
<br />
[9] Leon A. Gatys, Alexander S. Ecker, Matthias Bethge, A Neural Algorithm of Artistic Style, https://arxiv.org/abs/1508.06576<br />
<br />
[10] Karen Simonyan et al. Very Deep Convolutional Networks for Large-Scale Image Recognition<br />
<br />
[11] VGG Architectures - [http://www.robots.ox.ac.uk/~vgg/research/very_deep/| More Details]<br />
<br />
[12] Mechrez, R., Shechtman, E., & Zelnik-Manor, L. (2017). Photorealistic Style Transfer with Screened Poisson Equation. arXiv preprint arXiv:1709.09828.<br />
<br />
[13] Y. Li, C. Fang, J. Yang, Z. Wang, X. Lu, and M.-H. Yang. Diversified texture synthesis with feed-forward networks. In CVPR, 2017<br />
<br />
[14] D. Chen, L. Yuan, J. Liao, N. Yu, and G. Hua. Stylebank: An explicit representation for neural image style transfer. In CVPR, 2017<br />
<br />
Implementation Example: https://github.com/titu1994/Neural-Style-Transfer</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Alternative_Neural_Network:_Exploring_Contexts_As_Early_As_Possible_For_Action_Recognition&diff=30816Deep Alternative Neural Network: Exploring Contexts As Early As Possible For Action Recognition2017-11-19T18:36:05Z<p>Nkadiyar: /* References */</p>
<hr />
<div>==Introduction==<br />
<br />
Action recognition deals with recognizing and classifying the actions or activities done by humans or other agents in a video clip. In action recognition, Contexts contribute semantic clues for action recognition in video(See Fig below[8]). Conventional Neural Networks [1,2,3] and their shifted version 3D CNNs [4,5,6] have been employed in action recognition but they identify and aggregate the contexts at later stages. <br />
[[File:ActionRecognition1.jpg|center|400px|border|context and action region]]<br />
<br />
The authors have come up with a strategy to identify contexts in the videos as early as possible and leverage their evolutions for action recognition. Contexts contribute semantic clues for action recognition in videos. The network themselves involve a lot of layers and the first layer typically being receptive fields (RF) output only extra local features. As we go deeper into the layers the Receptive Fields expand and we start getting the contexts. The authors identified that increasing the number of layers will only cause additional burden in terms of handling the parameters and contexts could be obtained even in the earlier stages. The authors also cite the papers [9,10] that relate the CNNs and the visual systems of our brain, one remarkable difference being the abundant recurrent connections in our brain compared to the forward connections in the CNNs. <br />
<br />
The main contributions in the paper can be summarized as follows: <br />
* A Deep Alternative Neural Network (DANN) is proposed for action recognition. <br />
* DANN consists of alternative volumetric convolutional and recurrent layers. <br />
* An adaptive method to determine the temporal size of the video clip <br />
* A volumetric pyramid pooling layer to resize the output before fully connected layers.<br />
<br />
=== Optic Flow ===<br />
Optical flow or optic flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene.<br />
It can be used for affordance perception, the ability to discern possibilities for action within the environment.<br />
<br />
==Deep Alternative Neural Network:==<br />
===Adaptive Network Input===<br />
The input size of the video clip is generally determined empirically and various approaches have been taken in the past with a different number of frames. For instance, many previous papers suggested to used shorter intervals of between 1 to 16 frames. However, more recent work[9] recognized that human-based actions often “span tens or hundreds of frames” and longer intervals such as 60 frames will outperform the one with shorter interval. However, there’s still no systematic way of determining number of frames for input size of the network. This serves the motives for the authors of this paper to develop this adaptive method. Past research shows that motion energy intensity induced by human activity exhibits a regular periodicity. This signal can be approximately estimated by optical flow computation as shown in Figure 1, and is particularly suitable to address our temporal estimation due to: <br />
* the local minima and maxima landmarks probably correspond to characteristic gesture and motion <br />
* it is relatively robust to changes in camera viewpoint.<br />
<br />
The authors have come up with an adaptive method to automatically select the most discriminative video fragments using the density of optical flow energy which exhibits regular periodicity. According to Wikipedia, optical flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene, and optical flow methods try to calculate the motion between two image frames which are taken at different times. The optimal flow energy of an optical field $(v_{x},v_{y})$ is defined as follows <br />
<br />
:<math>e(I)=\underset{(x,y)\in\mathbb{P}}{\operatorname{\Sigma}} ||v_{x}(x,y),v_{y}(x,y)||_{2}</math><br />
<br />
Here, P is the pixel level set of selected interest points. They locate the local minima and maxima landmarks $\{t\}$ of $\epsilon = \{e(I_1),\dots,e(I_t)\}$ and for each two consecutive landmarks create a video fragment $s$ by extracting the frames $s = \{I_{t-1},\dots,I_t\}$.<br />
<br />
[[File:golfswing.png]]<br />
<br />
===Alternative Layer===<br />
This is a key layer consisting of a standard volumetric convolutional layer followed by a designed recurrent layer. Volumetric convolutional extracts features from local neighborhoods and a recurrent layer is applied to the output and it proceeds iteratively for T times. The input of a unit at position (x,y,z) in the jth feature map of the ith AL in time t is given by,<br />
<br />
:<math>u_{ij}^{xyz}(t) = u_{ij}^{xyz}(0) + f(w_{ij}^{r}u_{ij}^{xyz}(t-1)) + b_{ij} \\ <br />
u_{ij}^{xyz}(0) = f(w_{i-1}^{c}u_{(i-1)j}^{xyz}) <br />
</math><br />
<br />
U(0): feed forward output of volumetric convolutional layer. <br />
U(t-1) : recurrent input of previous time <br />
w_{k}^{c} and w_{k}^{r}: vectorized feed-forward kernels and recurrent kernels respectively <br />
f: ReLU function<br />
<br />
Figure 3 depicts this structure:<br />
[[File:unfolded.PNG|1000px]]<br />
<br />
===Volumetric Pyramid Pooling Layer===<br />
<br />
[[File:Volumetric Pyramid Pooling Layer.png|thumb|550px|Figure 2: Volumetric Pyramid Pooling Layer]]<br />
The authors have replaced the last pooling layer with a volumetric pyramid pooling layer as we need fixed-length vectors for the fully connected layers and the AL accepts video clips of arbitrary sizes and produces outputs of variable sizes. The authors have used the max pooling to pool the responses of each kernel in each volumetric bin. The outputs are kM dimensional vectors where:<br />
<br />
M: number of bins <br />
<br />
K: Number of kernels in the last alternative layer.<br />
<br />
This layer structure allows not only for arbitrary-length videos, but also arbitrary aspect ratios and scales.<br />
<br />
==Overall Architecture== <br />
[[File:DANN Architecture.png|thumb|550px|Figure 3:DANN Architecture]]<br />
The following are the components of the DANN (as shown in Figure 3)<br />
* 6 Alternative layers with 64, 128, 256, 256, 512 and 512 kernel response maps <br />
* 5 ReLU and volumetric pooling layers <br />
* 1 volumetric pyramid pooling layer <br />
* 3 fully connected layers of size 2048 each <br />
* A softmax layer<br />
<br />
==Implementation details==<br />
The authors have used the Torch toolbox platform for Implementations of volumetric convolutions, recurrent layers and optimizations. They have used a technique called as random clipping for data augmentation, in which they select a point randomly from the input video of fixed size 80x80xt after determining the temporal size t. For training the network the authors have used SGD applied to mini-batches of size 30 with a negative log likelihood criterion. Training is done by minimizing the cross-entropy loss function using backpropagation through time algorithm (BPTT). During testing, they applied a video clip divided into 80x80xt clips with a stride of 4 frames followed by testing with 10 crops. Final score is the average of all clip-level scores and the crop scores.<br />
Data augmentation techniques such as the multi-scale cropping method have been evaluated due to the recent success in the state-of-the-art performance displayed by Very Deep Two-stream ConvNets. Going by intuition, the corner cropping strategy could provide better results ( based on trade-off degree) since the receptive fields can focus harder on the central regions of the video frames [7].<br />
<br />
==Evaluations==<br />
===Datasets:===<br />
* The datasets used in the evaluation are UCF101 and HMDB51 <br />
* UCF101 – 13K videos annotated into 101 classes <br />
* HMDB51 – 6.8K videos with 51 actions. <br />
* Three training and test splits are provided <br />
* Performance measured by mean classification accuracy across the splits. <br />
* UCF101 split – 9.5K videos; HMDB51 – 3.7K training videos.<br />
<br />
===Quantitative Results===<br />
The authors used three types of optical flows, viz., sparse, RGB and TVL1 and found that TVL1 is suitable as action recognition is more easy to learn from motion information compared to raw pixel values. The influence of data augmentation is also studied. The baseline being sliding window with 75% overlap, the authors observed that the random clipping and multi-scale clipping outperformed the baseline on the UCF101 split 1 dataset. The authors were able to prove that the adaptive temporal length was able to give a boost of 4.2% when compared with architectures that had fixed-size temporal length. Experiments were also conducted to see if the learnings done in one dataset could improve the accuracy of another dataset. Fine tuning HMDB51 from UCF101 boosted the performance from 56.4% to 62.5%. The authors also observed that increasing the AL layers improves the performance as larger contexts are being embedded into the DANN. The DANN achieved an overall accuracy of 65.9% and 91.6% on HMDB51 and UCF101 respectively.<br />
<br />
<br />
[[File:Performance Comparison of different input modalities.png]]<br />
<br />
===Qualitative Analysis===<br />
The authors have discussed the quality of the prediction in the video clips taking examples of two different scenes involving bowling and haircut. In the bowling scene, the adaptive temporal choice used by DANN could aggregate more reasonable semantic structures and hence it leveraged reasonable video clips as input. On the other hand, the performance on the haircut video clip was not up to the mark as the rich contexts provided by the DANN was not helpful in a setting with simple actions performed in a simple background.<br />
<br />
==Conclusion==<br />
* Deep alternative neural network is introduced for action recognition. <br />
* DANN consists of volumetric convolutional layer and a recurrent layer. <br />
* The authors have experimented with different datasets like HMDB51 and UCF101 with different scenarios and compared the * * performance of DANN with other approaches. <br />
* The spatial size is still chosen in an ad hoc manner and this can be an area of improvement. <br />
* There are prospects for studying action tube which is a more compact input.<br />
<br />
Github code: https://github.com/wangjinzhuo/DANN<br />
<br />
==References==<br />
<br />
[1] Andrej Karpathy, George Toderici, Sachin Shetty, Tommy Leung, Rahul Sukthankar, and Li FeiFei. Large-scale video classification with convolutional neural networks. In CVPR, pages 1725–1732, 2014 <br />
<br />
[2] Karen Simonyan and Andrew Zisserman. Two-stream convolutional networks for action recognition in videos. In NIPS, pages 568–576, 2014. <br />
<br />
[3]Limin Wang, Yu Qiao, and Xiaoou Tang. Action recognition with trajectory-pooled deepconvolutional descriptors. In CVPR, pages 4305–4314, 2015. <br />
<br />
[4] Shuiwang Ji, Wei Xu, Ming Yang, and Kai Yu. 3d convolutional neural networks for human action recognition. TPAMI, 35(1):221–231, 2013. <br />
<br />
[5] Du Tran, Lubomir Bourdev, Rob Fergus, Lorenzo Torresani, and Manohar Paluri. Learning spatiotemporal features with 3d convolutional networks. In ICCV, pages 4489–4497, 2015. <br />
<br />
[6]Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action recognition. arXiv preprint arXiv:1604.04494, 2016. <br />
<br />
[7]Limin Wang, Yuanjun Xiong, Zhe Wang, Yu Qiao. Towards Good Practices for Very Deep Two-Stream ConvNets. arXiv preprint arXiv:1507.02159 , 2015. <br />
<br />
[8] IEEE International Symposium on Multimedia 2013 <br />
<br />
[9] Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action<br />
recognition. arXiv preprint arXiv:1604.04494, 2016<br />
<br />
[10] https://en.wikipedia.org/wiki/Optical_flow<br />
<br />
[36] Christopher Zach, Thomas Pock, and Horst Bischof. A duality based approach for realtime tv-l<br />
1 optical flow. In Pattern Recognition, pages 214–223. 2007.<br />
<br />
A list of expert reviews: http://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/480.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Alternative_Neural_Network:_Exploring_Contexts_As_Early_As_Possible_For_Action_Recognition&diff=30815Deep Alternative Neural Network: Exploring Contexts As Early As Possible For Action Recognition2017-11-19T18:35:30Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>==Introduction==<br />
<br />
Action recognition deals with recognizing and classifying the actions or activities done by humans or other agents in a video clip. In action recognition, Contexts contribute semantic clues for action recognition in video(See Fig below[8]). Conventional Neural Networks [1,2,3] and their shifted version 3D CNNs [4,5,6] have been employed in action recognition but they identify and aggregate the contexts at later stages. <br />
[[File:ActionRecognition1.jpg|center|400px|border|context and action region]]<br />
<br />
The authors have come up with a strategy to identify contexts in the videos as early as possible and leverage their evolutions for action recognition. Contexts contribute semantic clues for action recognition in videos. The network themselves involve a lot of layers and the first layer typically being receptive fields (RF) output only extra local features. As we go deeper into the layers the Receptive Fields expand and we start getting the contexts. The authors identified that increasing the number of layers will only cause additional burden in terms of handling the parameters and contexts could be obtained even in the earlier stages. The authors also cite the papers [9,10] that relate the CNNs and the visual systems of our brain, one remarkable difference being the abundant recurrent connections in our brain compared to the forward connections in the CNNs. <br />
<br />
The main contributions in the paper can be summarized as follows: <br />
* A Deep Alternative Neural Network (DANN) is proposed for action recognition. <br />
* DANN consists of alternative volumetric convolutional and recurrent layers. <br />
* An adaptive method to determine the temporal size of the video clip <br />
* A volumetric pyramid pooling layer to resize the output before fully connected layers.<br />
<br />
=== Optic Flow ===<br />
Optical flow or optic flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene.<br />
It can be used for affordance perception, the ability to discern possibilities for action within the environment.<br />
<br />
==Deep Alternative Neural Network:==<br />
===Adaptive Network Input===<br />
The input size of the video clip is generally determined empirically and various approaches have been taken in the past with a different number of frames. For instance, many previous papers suggested to used shorter intervals of between 1 to 16 frames. However, more recent work[9] recognized that human-based actions often “span tens or hundreds of frames” and longer intervals such as 60 frames will outperform the one with shorter interval. However, there’s still no systematic way of determining number of frames for input size of the network. This serves the motives for the authors of this paper to develop this adaptive method. Past research shows that motion energy intensity induced by human activity exhibits a regular periodicity. This signal can be approximately estimated by optical flow computation as shown in Figure 1, and is particularly suitable to address our temporal estimation due to: <br />
* the local minima and maxima landmarks probably correspond to characteristic gesture and motion <br />
* it is relatively robust to changes in camera viewpoint.<br />
<br />
The authors have come up with an adaptive method to automatically select the most discriminative video fragments using the density of optical flow energy which exhibits regular periodicity. According to Wikipedia, optical flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene, and optical flow methods try to calculate the motion between two image frames which are taken at different times. The optimal flow energy of an optical field $(v_{x},v_{y})$ is defined as follows <br />
<br />
:<math>e(I)=\underset{(x,y)\in\mathbb{P}}{\operatorname{\Sigma}} ||v_{x}(x,y),v_{y}(x,y)||_{2}</math><br />
<br />
Here, P is the pixel level set of selected interest points. They locate the local minima and maxima landmarks $\{t\}$ of $\epsilon = \{e(I_1),\dots,e(I_t)\}$ and for each two consecutive landmarks create a video fragment $s$ by extracting the frames $s = \{I_{t-1},\dots,I_t\}$.<br />
<br />
[[File:golfswing.png]]<br />
<br />
===Alternative Layer===<br />
This is a key layer consisting of a standard volumetric convolutional layer followed by a designed recurrent layer. Volumetric convolutional extracts features from local neighborhoods and a recurrent layer is applied to the output and it proceeds iteratively for T times. The input of a unit at position (x,y,z) in the jth feature map of the ith AL in time t is given by,<br />
<br />
:<math>u_{ij}^{xyz}(t) = u_{ij}^{xyz}(0) + f(w_{ij}^{r}u_{ij}^{xyz}(t-1)) + b_{ij} \\ <br />
u_{ij}^{xyz}(0) = f(w_{i-1}^{c}u_{(i-1)j}^{xyz}) <br />
</math><br />
<br />
U(0): feed forward output of volumetric convolutional layer. <br />
U(t-1) : recurrent input of previous time <br />
w_{k}^{c} and w_{k}^{r}: vectorized feed-forward kernels and recurrent kernels respectively <br />
f: ReLU function<br />
<br />
Figure 3 depicts this structure:<br />
[[File:unfolded.PNG|1000px]]<br />
<br />
===Volumetric Pyramid Pooling Layer===<br />
<br />
[[File:Volumetric Pyramid Pooling Layer.png|thumb|550px|Figure 2: Volumetric Pyramid Pooling Layer]]<br />
The authors have replaced the last pooling layer with a volumetric pyramid pooling layer as we need fixed-length vectors for the fully connected layers and the AL accepts video clips of arbitrary sizes and produces outputs of variable sizes. The authors have used the max pooling to pool the responses of each kernel in each volumetric bin. The outputs are kM dimensional vectors where:<br />
<br />
M: number of bins <br />
<br />
K: Number of kernels in the last alternative layer.<br />
<br />
This layer structure allows not only for arbitrary-length videos, but also arbitrary aspect ratios and scales.<br />
<br />
==Overall Architecture== <br />
[[File:DANN Architecture.png|thumb|550px|Figure 3:DANN Architecture]]<br />
The following are the components of the DANN (as shown in Figure 3)<br />
* 6 Alternative layers with 64, 128, 256, 256, 512 and 512 kernel response maps <br />
* 5 ReLU and volumetric pooling layers <br />
* 1 volumetric pyramid pooling layer <br />
* 3 fully connected layers of size 2048 each <br />
* A softmax layer<br />
<br />
==Implementation details==<br />
The authors have used the Torch toolbox platform for Implementations of volumetric convolutions, recurrent layers and optimizations. They have used a technique called as random clipping for data augmentation, in which they select a point randomly from the input video of fixed size 80x80xt after determining the temporal size t. For training the network the authors have used SGD applied to mini-batches of size 30 with a negative log likelihood criterion. Training is done by minimizing the cross-entropy loss function using backpropagation through time algorithm (BPTT). During testing, they applied a video clip divided into 80x80xt clips with a stride of 4 frames followed by testing with 10 crops. Final score is the average of all clip-level scores and the crop scores.<br />
Data augmentation techniques such as the multi-scale cropping method have been evaluated due to the recent success in the state-of-the-art performance displayed by Very Deep Two-stream ConvNets. Going by intuition, the corner cropping strategy could provide better results ( based on trade-off degree) since the receptive fields can focus harder on the central regions of the video frames [7].<br />
<br />
==Evaluations==<br />
===Datasets:===<br />
* The datasets used in the evaluation are UCF101 and HMDB51 <br />
* UCF101 – 13K videos annotated into 101 classes <br />
* HMDB51 – 6.8K videos with 51 actions. <br />
* Three training and test splits are provided <br />
* Performance measured by mean classification accuracy across the splits. <br />
* UCF101 split – 9.5K videos; HMDB51 – 3.7K training videos.<br />
<br />
===Quantitative Results===<br />
The authors used three types of optical flows, viz., sparse, RGB and TVL1 and found that TVL1 is suitable as action recognition is more easy to learn from motion information compared to raw pixel values. The influence of data augmentation is also studied. The baseline being sliding window with 75% overlap, the authors observed that the random clipping and multi-scale clipping outperformed the baseline on the UCF101 split 1 dataset. The authors were able to prove that the adaptive temporal length was able to give a boost of 4.2% when compared with architectures that had fixed-size temporal length. Experiments were also conducted to see if the learnings done in one dataset could improve the accuracy of another dataset. Fine tuning HMDB51 from UCF101 boosted the performance from 56.4% to 62.5%. The authors also observed that increasing the AL layers improves the performance as larger contexts are being embedded into the DANN. The DANN achieved an overall accuracy of 65.9% and 91.6% on HMDB51 and UCF101 respectively.<br />
<br />
<br />
[[File:Performance Comparison of different input modalities.png]]<br />
<br />
===Qualitative Analysis===<br />
The authors have discussed the quality of the prediction in the video clips taking examples of two different scenes involving bowling and haircut. In the bowling scene, the adaptive temporal choice used by DANN could aggregate more reasonable semantic structures and hence it leveraged reasonable video clips as input. On the other hand, the performance on the haircut video clip was not up to the mark as the rich contexts provided by the DANN was not helpful in a setting with simple actions performed in a simple background.<br />
<br />
==Conclusion==<br />
* Deep alternative neural network is introduced for action recognition. <br />
* DANN consists of volumetric convolutional layer and a recurrent layer. <br />
* The authors have experimented with different datasets like HMDB51 and UCF101 with different scenarios and compared the * * performance of DANN with other approaches. <br />
* The spatial size is still chosen in an ad hoc manner and this can be an area of improvement. <br />
* There are prospects for studying action tube which is a more compact input.<br />
<br />
Github code: https://github.com/wangjinzhuo/DANN<br />
<br />
==References==<br />
<br />
[1] Andrej Karpathy, George Toderici, Sachin Shetty, Tommy Leung, Rahul Sukthankar, and Li FeiFei. Large-scale video classification with convolutional neural networks. In CVPR, pages 1725–1732, 2014 <br />
<br />
[2] Karen Simonyan and Andrew Zisserman. Two-stream convolutional networks for action recognition in videos. In NIPS, pages 568–576, 2014. <br />
<br />
[3]Limin Wang, Yu Qiao, and Xiaoou Tang. Action recognition with trajectory-pooled deepconvolutional descriptors. In CVPR, pages 4305–4314, 2015. <br />
<br />
[4] Shuiwang Ji, Wei Xu, Ming Yang, and Kai Yu. 3d convolutional neural networks for human action recognition. TPAMI, 35(1):221–231, 2013. <br />
<br />
[5] Du Tran, Lubomir Bourdev, Rob Fergus, Lorenzo Torresani, and Manohar Paluri. Learning spatiotemporal features with 3d convolutional networks. In ICCV, pages 4489–4497, 2015. <br />
<br />
[6]Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action recognition. arXiv preprint arXiv:1604.04494, 2016. <br />
<br />
[7]Limin Wang, Yuanjun Xiong, Zhe Wang, Yu Qiao. Towards Good Practices for Very Deep Two-Stream ConvNets. arXiv preprint arXiv:1507.02159 , 2015. <br />
<br />
[8] IEEE International Symposium on Multimedia 2013 <br />
<br />
[9] Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action<br />
recognition. arXiv preprint arXiv:1604.04494, 2016<br />
<br />
[36] Christopher Zach, Thomas Pock, and Horst Bischof. A duality based approach for realtime tv-l<br />
1 optical flow. In Pattern Recognition, pages 214–223. 2007.<br />
<br />
A list of expert reviews: http://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/480.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Alternative_Neural_Network:_Exploring_Contexts_As_Early_As_Possible_For_Action_Recognition&diff=30813Deep Alternative Neural Network: Exploring Contexts As Early As Possible For Action Recognition2017-11-19T18:27:23Z<p>Nkadiyar: /* Adaptive Network Input */</p>
<hr />
<div>==Introduction==<br />
<br />
Action recognition deals with recognizing and classifying the actions or activities done by humans or other agents in a video clip. In action recognition, Contexts contribute semantic clues for action recognition in video(See Fig below[8]). Conventional Neural Networks [1,2,3] and their shifted version 3D CNNs [4,5,6] have been employed in action recognition but they identify and aggregate the contexts at later stages. <br />
[[File:ActionRecognition1.jpg|center|400px|border|context and action region]]<br />
<br />
The authors have come up with a strategy to identify contexts in the videos as early as possible and leverage their evolutions for action recognition. Contexts contribute semantic clues for action recognition in videos. The network themselves involve a lot of layers and the first layer typically being receptive fields (RF) output only extra local features. As we go deeper into the layers the Receptive Fields expand and we start getting the contexts. The authors identified that increasing the number of layers will only cause additional burden in terms of handling the parameters and contexts could be obtained even in the earlier stages. The authors also cite the papers [9,10] that relate the CNNs and the visual systems of our brain, one remarkable difference being the abundant recurrent connections in our brain compared to the forward connections in the CNNs. <br />
<br />
The main contributions in the paper can be summarized as follows: <br />
* A Deep Alternative Neural Network (DANN) is proposed for action recognition. <br />
* DANN consists of alternative volumetric convolutional and recurrent layers. <br />
* An adaptive method to determine the temporal size of the video clip <br />
* A volumetric pyramid pooling layer to resize the output before fully connected layers.<br />
<br />
==Deep Alternative Neural Network:==<br />
===Adaptive Network Input===<br />
The input size of the video clip is generally determined empirically and various approaches have been taken in the past with a different number of frames. For instance, many previous papers suggested to used shorter intervals of between 1 to 16 frames. However, more recent work[9] recognized that human-based actions often “span tens or hundreds of frames” and longer intervals such as 60 frames will outperform the one with shorter interval. However, there’s still no systematic way of determining number of frames for input size of the network. This serves the motives for the authors of this paper to develop this adaptive method. Past research shows that motion energy intensity induced by human activity exhibits a regular periodicity. This signal can be approximately estimated by optical flow computation as shown in Figure 1, and is particularly suitable to address our temporal estimation due to: <br />
* the local minima and maxima landmarks probably correspond to characteristic gesture and motion <br />
* it is relatively robust to changes in camera viewpoint.<br />
<br />
The authors have come up with an adaptive method to automatically select the most discriminative video fragments using the density of optical flow energy which exhibits regular periodicity. According to Wikipedia, optical flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene, and optical flow methods try to calculate the motion between two image frames which are taken at different times. The optimal flow energy of an optical field $(v_{x},v_{y})$ is defined as follows <br />
<br />
:<math>e(I)=\underset{(x,y)\in\mathbb{P}}{\operatorname{\Sigma}} ||v_{x}(x,y),v_{y}(x,y)||_{2}</math><br />
<br />
Here, P is the pixel level set of selected interest points. They locate the local minima and maxima landmarks $\{t\}$ of $\epsilon = \{e(I_1),\dots,e(I_t)\}$ and for each two consecutive landmarks create a video fragment $s$ by extracting the frames $s = \{I_{t-1},\dots,I_t\}$.<br />
<br />
[[File:golfswing.png]]<br />
<br />
===Alternative Layer===<br />
This is a key layer consisting of a standard volumetric convolutional layer followed by a designed recurrent layer. Volumetric convolutional extracts features from local neighborhoods and a recurrent layer is applied to the output and it proceeds iteratively for T times. The input of a unit at position (x,y,z) in the jth feature map of the ith AL in time t is given by,<br />
<br />
:<math>u_{ij}^{xyz}(t) = u_{ij}^{xyz}(0) + f(w_{ij}^{r}u_{ij}^{xyz}(t-1)) + b_{ij} \\ <br />
u_{ij}^{xyz}(0) = f(w_{i-1}^{c}u_{(i-1)j}^{xyz}) <br />
</math><br />
<br />
U(0): feed forward output of volumetric convolutional layer. <br />
U(t-1) : recurrent input of previous time <br />
w_{k}^{c} and w_{k}^{r}: vectorized feed-forward kernels and recurrent kernels respectively <br />
f: ReLU function<br />
<br />
Figure 3 depicts this structure:<br />
[[File:unfolded.PNG|1000px]]<br />
<br />
===Volumetric Pyramid Pooling Layer===<br />
<br />
[[File:Volumetric Pyramid Pooling Layer.png|thumb|550px|Figure 2: Volumetric Pyramid Pooling Layer]]<br />
The authors have replaced the last pooling layer with a volumetric pyramid pooling layer as we need fixed-length vectors for the fully connected layers and the AL accepts video clips of arbitrary sizes and produces outputs of variable sizes. The authors have used the max pooling to pool the responses of each kernel in each volumetric bin. The outputs are kM dimensional vectors where:<br />
<br />
M: number of bins <br />
<br />
K: Number of kernels in the last alternative layer.<br />
<br />
This layer structure allows not only for arbitrary-length videos, but also arbitrary aspect ratios and scales.<br />
<br />
==Overall Architecture== <br />
[[File:DANN Architecture.png|thumb|550px|Figure 3:DANN Architecture]]<br />
The following are the components of the DANN (as shown in Figure 3)<br />
* 6 Alternative layers with 64, 128, 256, 256, 512 and 512 kernel response maps <br />
* 5 ReLU and volumetric pooling layers <br />
* 1 volumetric pyramid pooling layer <br />
* 3 fully connected layers of size 2048 each <br />
* A softmax layer<br />
<br />
==Implementation details==<br />
The authors have used the Torch toolbox platform for Implementations of volumetric convolutions, recurrent layers and optimizations. They have used a technique called as random clipping for data augmentation, in which they select a point randomly from the input video of fixed size 80x80xt after determining the temporal size t. For training the network the authors have used SGD applied to mini-batches of size 30 with a negative log likelihood criterion. Training is done by minimizing the cross-entropy loss function using backpropagation through time algorithm (BPTT). During testing, they applied a video clip divided into 80x80xt clips with a stride of 4 frames followed by testing with 10 crops. Final score is the average of all clip-level scores and the crop scores.<br />
Data augmentation techniques such as the multi-scale cropping method have been evaluated due to the recent success in the state-of-the-art performance displayed by Very Deep Two-stream ConvNets. Going by intuition, the corner cropping strategy could provide better results ( based on trade-off degree) since the receptive fields can focus harder on the central regions of the video frames [7].<br />
<br />
==Evaluations==<br />
===Datasets:===<br />
* The datasets used in the evaluation are UCF101 and HMDB51 <br />
* UCF101 – 13K videos annotated into 101 classes <br />
* HMDB51 – 6.8K videos with 51 actions. <br />
* Three training and test splits are provided <br />
* Performance measured by mean classification accuracy across the splits. <br />
* UCF101 split – 9.5K videos; HMDB51 – 3.7K training videos.<br />
<br />
===Quantitative Results===<br />
The authors used three types of optical flows, viz., sparse, RGB and TVL1 and found that TVL1 is suitable as action recognition is more easy to learn from motion information compared to raw pixel values. The influence of data augmentation is also studied. The baseline being sliding window with 75% overlap, the authors observed that the random clipping and multi-scale clipping outperformed the baseline on the UCF101 split 1 dataset. The authors were able to prove that the adaptive temporal length was able to give a boost of 4.2% when compared with architectures that had fixed-size temporal length. Experiments were also conducted to see if the learnings done in one dataset could improve the accuracy of another dataset. Fine tuning HMDB51 from UCF101 boosted the performance from 56.4% to 62.5%. The authors also observed that increasing the AL layers improves the performance as larger contexts are being embedded into the DANN. The DANN achieved an overall accuracy of 65.9% and 91.6% on HMDB51 and UCF101 respectively.<br />
<br />
<br />
[[File:Performance Comparison of different input modalities.png]]<br />
<br />
===Qualitative Analysis===<br />
The authors have discussed the quality of the prediction in the video clips taking examples of two different scenes involving bowling and haircut. In the bowling scene, the adaptive temporal choice used by DANN could aggregate more reasonable semantic structures and hence it leveraged reasonable video clips as input. On the other hand, the performance on the haircut video clip was not up to the mark as the rich contexts provided by the DANN was not helpful in a setting with simple actions performed in a simple background.<br />
<br />
==Conclusion==<br />
* Deep alternative neural network is introduced for action recognition. <br />
* DANN consists of volumetric convolutional layer and a recurrent layer. <br />
* The authors have experimented with different datasets like HMDB51 and UCF101 with different scenarios and compared the * * performance of DANN with other approaches. <br />
* The spatial size is still chosen in an ad hoc manner and this can be an area of improvement. <br />
* There are prospects for studying action tube which is a more compact input.<br />
<br />
Github code: https://github.com/wangjinzhuo/DANN<br />
<br />
==References==<br />
<br />
[1] Andrej Karpathy, George Toderici, Sachin Shetty, Tommy Leung, Rahul Sukthankar, and Li FeiFei. Large-scale video classification with convolutional neural networks. In CVPR, pages 1725–1732, 2014 <br />
<br />
[2] Karen Simonyan and Andrew Zisserman. Two-stream convolutional networks for action recognition in videos. In NIPS, pages 568–576, 2014. <br />
<br />
[3]Limin Wang, Yu Qiao, and Xiaoou Tang. Action recognition with trajectory-pooled deepconvolutional descriptors. In CVPR, pages 4305–4314, 2015. <br />
<br />
[4] Shuiwang Ji, Wei Xu, Ming Yang, and Kai Yu. 3d convolutional neural networks for human action recognition. TPAMI, 35(1):221–231, 2013. <br />
<br />
[5] Du Tran, Lubomir Bourdev, Rob Fergus, Lorenzo Torresani, and Manohar Paluri. Learning spatiotemporal features with 3d convolutional networks. In ICCV, pages 4489–4497, 2015. <br />
<br />
[6]Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action recognition. arXiv preprint arXiv:1604.04494, 2016. <br />
<br />
[7]Limin Wang, Yuanjun Xiong, Zhe Wang, Yu Qiao. Towards Good Practices for Very Deep Two-Stream ConvNets. arXiv preprint arXiv:1507.02159 , 2015. <br />
<br />
[8] IEEE International Symposium on Multimedia 2013 <br />
<br />
[9] Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action<br />
recognition. arXiv preprint arXiv:1604.04494, 2016<br />
<br />
[36] Christopher Zach, Thomas Pock, and Horst Bischof. A duality based approach for realtime tv-l<br />
1 optical flow. In Pattern Recognition, pages 214–223. 2007.<br />
<br />
A list of expert reviews: http://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/480.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Alternative_Neural_Network:_Exploring_Contexts_As_Early_As_Possible_For_Action_Recognition&diff=30812Deep Alternative Neural Network: Exploring Contexts As Early As Possible For Action Recognition2017-11-19T18:26:17Z<p>Nkadiyar: </p>
<hr />
<div>==Introduction==<br />
<br />
Action recognition deals with recognizing and classifying the actions or activities done by humans or other agents in a video clip. In action recognition, Contexts contribute semantic clues for action recognition in video(See Fig below[8]). Conventional Neural Networks [1,2,3] and their shifted version 3D CNNs [4,5,6] have been employed in action recognition but they identify and aggregate the contexts at later stages. <br />
[[File:ActionRecognition1.jpg|center|400px|border|context and action region]]<br />
<br />
The authors have come up with a strategy to identify contexts in the videos as early as possible and leverage their evolutions for action recognition. Contexts contribute semantic clues for action recognition in videos. The network themselves involve a lot of layers and the first layer typically being receptive fields (RF) output only extra local features. As we go deeper into the layers the Receptive Fields expand and we start getting the contexts. The authors identified that increasing the number of layers will only cause additional burden in terms of handling the parameters and contexts could be obtained even in the earlier stages. The authors also cite the papers [9,10] that relate the CNNs and the visual systems of our brain, one remarkable difference being the abundant recurrent connections in our brain compared to the forward connections in the CNNs. <br />
<br />
The main contributions in the paper can be summarized as follows: <br />
* A Deep Alternative Neural Network (DANN) is proposed for action recognition. <br />
* DANN consists of alternative volumetric convolutional and recurrent layers. <br />
* An adaptive method to determine the temporal size of the video clip <br />
* A volumetric pyramid pooling layer to resize the output before fully connected layers.<br />
<br />
==Deep Alternative Neural Network:==<br />
===Adaptive Network Input===<br />
The input size of the video clip is generally determined empirically and various approaches have been taken in the past with a different number of frames. For instance, many previous papers suggested to used shorter intervals of between 1 to 16 frames. However, more recent work[9] recognized that human-based actions often “span tens or hundreds of frames” and longer intervals such as 60 frames will outperform the one with shorter interval. However, there’s still no systematic way of determining number of frames for input size of the network. This serves the motives for the authors of this paper to develop this adaptive method. Past research shows that motion energy intensity induced by human activity exhibits has a regular periodicity. This signal can be approximately estimated by optical flow computation as shown in Figure 1, and is particularly suitable to address our temporal estimation due to: <br />
* the local minima and maxima landmarks probably correspond to characteristic gesture and motion <br />
* it is relatively robust to changes in camera viewpoint.<br />
<br />
The authors have come up with an adaptive method to automatically select the most discriminative video fragments using the density of optical flow energy which exhibits regular periodicity. According to Wikipedia, optical flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene, and optical flow methods try to calculate the motion between two image frames which are taken at different times. The optimal flow energy of an optical field $(v_{x},v_{y})$ is defined as follows <br />
<br />
:<math>e(I)=\underset{(x,y)\in\mathbb{P}}{\operatorname{\Sigma}} ||v_{x}(x,y),v_{y}(x,y)||_{2}</math><br />
<br />
Here, P is the pixel level set of selected interest points. They locate the local minima and maxima landmarks $\{t\}$ of $\epsilon = \{e(I_1),\dots,e(I_t)\}$ and for each two consecutive landmarks create a video fragment $s$ by extracting the frames $s = \{I_{t-1},\dots,I_t\}$.<br />
<br />
[[File:golfswing.png]]<br />
<br />
===Alternative Layer===<br />
This is a key layer consisting of a standard volumetric convolutional layer followed by a designed recurrent layer. Volumetric convolutional extracts features from local neighborhoods and a recurrent layer is applied to the output and it proceeds iteratively for T times. The input of a unit at position (x,y,z) in the jth feature map of the ith AL in time t is given by,<br />
<br />
:<math>u_{ij}^{xyz}(t) = u_{ij}^{xyz}(0) + f(w_{ij}^{r}u_{ij}^{xyz}(t-1)) + b_{ij} \\ <br />
u_{ij}^{xyz}(0) = f(w_{i-1}^{c}u_{(i-1)j}^{xyz}) <br />
</math><br />
<br />
U(0): feed forward output of volumetric convolutional layer. <br />
U(t-1) : recurrent input of previous time <br />
w_{k}^{c} and w_{k}^{r}: vectorized feed-forward kernels and recurrent kernels respectively <br />
f: ReLU function<br />
<br />
Figure 3 depicts this structure:<br />
[[File:unfolded.PNG|1000px]]<br />
<br />
===Volumetric Pyramid Pooling Layer===<br />
<br />
[[File:Volumetric Pyramid Pooling Layer.png|thumb|550px|Figure 2: Volumetric Pyramid Pooling Layer]]<br />
The authors have replaced the last pooling layer with a volumetric pyramid pooling layer as we need fixed-length vectors for the fully connected layers and the AL accepts video clips of arbitrary sizes and produces outputs of variable sizes. The authors have used the max pooling to pool the responses of each kernel in each volumetric bin. The outputs are kM dimensional vectors where:<br />
<br />
M: number of bins <br />
<br />
K: Number of kernels in the last alternative layer.<br />
<br />
This layer structure allows not only for arbitrary-length videos, but also arbitrary aspect ratios and scales.<br />
<br />
==Overall Architecture== <br />
[[File:DANN Architecture.png|thumb|550px|Figure 3:DANN Architecture]]<br />
The following are the components of the DANN (as shown in Figure 3)<br />
* 6 Alternative layers with 64, 128, 256, 256, 512 and 512 kernel response maps <br />
* 5 ReLU and volumetric pooling layers <br />
* 1 volumetric pyramid pooling layer <br />
* 3 fully connected layers of size 2048 each <br />
* A softmax layer<br />
<br />
==Implementation details==<br />
The authors have used the Torch toolbox platform for Implementations of volumetric convolutions, recurrent layers and optimizations. They have used a technique called as random clipping for data augmentation, in which they select a point randomly from the input video of fixed size 80x80xt after determining the temporal size t. For training the network the authors have used SGD applied to mini-batches of size 30 with a negative log likelihood criterion. Training is done by minimizing the cross-entropy loss function using backpropagation through time algorithm (BPTT). During testing, they applied a video clip divided into 80x80xt clips with a stride of 4 frames followed by testing with 10 crops. Final score is the average of all clip-level scores and the crop scores.<br />
Data augmentation techniques such as the multi-scale cropping method have been evaluated due to the recent success in the state-of-the-art performance displayed by Very Deep Two-stream ConvNets. Going by intuition, the corner cropping strategy could provide better results ( based on trade-off degree) since the receptive fields can focus harder on the central regions of the video frames [7].<br />
<br />
==Evaluations==<br />
===Datasets:===<br />
* The datasets used in the evaluation are UCF101 and HMDB51 <br />
* UCF101 – 13K videos annotated into 101 classes <br />
* HMDB51 – 6.8K videos with 51 actions. <br />
* Three training and test splits are provided <br />
* Performance measured by mean classification accuracy across the splits. <br />
* UCF101 split – 9.5K videos; HMDB51 – 3.7K training videos.<br />
<br />
===Quantitative Results===<br />
The authors used three types of optical flows, viz., sparse, RGB and TVL1 and found that TVL1 is suitable as action recognition is more easy to learn from motion information compared to raw pixel values. The influence of data augmentation is also studied. The baseline being sliding window with 75% overlap, the authors observed that the random clipping and multi-scale clipping outperformed the baseline on the UCF101 split 1 dataset. The authors were able to prove that the adaptive temporal length was able to give a boost of 4.2% when compared with architectures that had fixed-size temporal length. Experiments were also conducted to see if the learnings done in one dataset could improve the accuracy of another dataset. Fine tuning HMDB51 from UCF101 boosted the performance from 56.4% to 62.5%. The authors also observed that increasing the AL layers improves the performance as larger contexts are being embedded into the DANN. The DANN achieved an overall accuracy of 65.9% and 91.6% on HMDB51 and UCF101 respectively.<br />
<br />
<br />
[[File:Performance Comparison of different input modalities.png]]<br />
<br />
===Qualitative Analysis===<br />
The authors have discussed the quality of the prediction in the video clips taking examples of two different scenes involving bowling and haircut. In the bowling scene, the adaptive temporal choice used by DANN could aggregate more reasonable semantic structures and hence it leveraged reasonable video clips as input. On the other hand, the performance on the haircut video clip was not up to the mark as the rich contexts provided by the DANN was not helpful in a setting with simple actions performed in a simple background.<br />
<br />
==Conclusion==<br />
* Deep alternative neural network is introduced for action recognition. <br />
* DANN consists of volumetric convolutional layer and a recurrent layer. <br />
* The authors have experimented with different datasets like HMDB51 and UCF101 with different scenarios and compared the * * performance of DANN with other approaches. <br />
* The spatial size is still chosen in an ad hoc manner and this can be an area of improvement. <br />
* There are prospects for studying action tube which is a more compact input.<br />
<br />
Github code: https://github.com/wangjinzhuo/DANN<br />
<br />
==References==<br />
<br />
[1] Andrej Karpathy, George Toderici, Sachin Shetty, Tommy Leung, Rahul Sukthankar, and Li FeiFei. Large-scale video classification with convolutional neural networks. In CVPR, pages 1725–1732, 2014 <br />
<br />
[2] Karen Simonyan and Andrew Zisserman. Two-stream convolutional networks for action recognition in videos. In NIPS, pages 568–576, 2014. <br />
<br />
[3]Limin Wang, Yu Qiao, and Xiaoou Tang. Action recognition with trajectory-pooled deepconvolutional descriptors. In CVPR, pages 4305–4314, 2015. <br />
<br />
[4] Shuiwang Ji, Wei Xu, Ming Yang, and Kai Yu. 3d convolutional neural networks for human action recognition. TPAMI, 35(1):221–231, 2013. <br />
<br />
[5] Du Tran, Lubomir Bourdev, Rob Fergus, Lorenzo Torresani, and Manohar Paluri. Learning spatiotemporal features with 3d convolutional networks. In ICCV, pages 4489–4497, 2015. <br />
<br />
[6]Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action recognition. arXiv preprint arXiv:1604.04494, 2016. <br />
<br />
[7]Limin Wang, Yuanjun Xiong, Zhe Wang, Yu Qiao. Towards Good Practices for Very Deep Two-Stream ConvNets. arXiv preprint arXiv:1507.02159 , 2015. <br />
<br />
[8] IEEE International Symposium on Multimedia 2013 <br />
<br />
[9] Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action<br />
recognition. arXiv preprint arXiv:1604.04494, 2016<br />
<br />
[36] Christopher Zach, Thomas Pock, and Horst Bischof. A duality based approach for realtime tv-l<br />
1 optical flow. In Pattern Recognition, pages 214–223. 2007.<br />
<br />
A list of expert reviews: http://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/480.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Alternative_Neural_Network:_Exploring_Contexts_As_Early_As_Possible_For_Action_Recognition&diff=30810Deep Alternative Neural Network: Exploring Contexts As Early As Possible For Action Recognition2017-11-19T18:23:01Z<p>Nkadiyar: /* References */</p>
<hr />
<div>==Introduction==<br />
<br />
Action recognition deals with recognizing and classifying the actions or activities done by humans or other agents in a video clip. In action recognition, Contexts contribute semantic clues for action recognition in video(See Fig below[8]). Conventional Neural Networks [1,2,3] and their shifted version 3D CNNs [4,5,6] have been employed in action recognition but they identify and aggregate the contexts at later stages. <br />
[[File:ActionRecognition1.jpg|center|400px|border|context and action region]]<br />
<br />
The authors have come up with a strategy to identify contexts in the videos as early as possible and leverage their evolutions for action recognition. Contexts contribute semantic clues for action recognition in videos. The network themselves involve a lot of layers and the first layer typically being receptive fields (RF) output only extra local features. As we go deeper into the layers the Receptive Fields expand and we start getting the contexts. The authors identified that increasing the number of layers will only cause additional burden in terms of handling the parameters and contexts could be obtained even in the earlier stages. The authors also cite the papers [9,10] that relate the CNNs and the visual systems of our brain, one remarkable difference being the abundant recurrent connections in our brain compared to the forward connections in the CNNs. <br />
<br />
The main contributions in the paper can be summarized as follows: <br />
* A Deep Alternative Neural Network (DANN) is proposed for action recognition. <br />
* DANN consists of alternative volumetric convolutional and recurrent layers. <br />
* An adaptive method to determine the temporal size of the video clip <br />
* A volumetric pyramid pooling layer to resize the output before fully connected layers.<br />
<br />
==Deep Alternative Neural Network:==<br />
===Adaptive Network Input===<br />
The input size of the video clip is generally determined empirically and various approaches have been taken in the past with a different number of frames. For instance, many previous papers suggested to used shorter intervals of between 1 to 16 frames. However, more recent work recognized that human-based actions often “span tens or hundreds of frames” and longer intervals such as 60 frames will outperform the one with shorter interval. However, there’s still no systematic way of determining number of frames for input size of the network. This serves the motives for the authors of this paper to develop this adaptive method. Past research shows that motion energy intensity induced by human activity exhibits has a regular periodicity. This signal can be approximately estimated by optical flow computation as shown in Figure 1, and is particularly suitable to address our temporal estimation due to: <br />
* the local minima and maxima landmarks probably correspond to characteristic gesture and motion <br />
* it is relatively robust to changes in camera viewpoint.<br />
<br />
The authors have come up with an adaptive method to automatically select the most discriminative video fragments using the density of optical flow energy which exhibits regular periodicity. According to Wikipedia, optical flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene, and optical flow methods try to calculate the motion between two image frames which are taken at different times. The optimal flow energy of an optical field $(v_{x},v_{y})$ is defined as follows <br />
<br />
:<math>e(I)=\underset{(x,y)\in\mathbb{P}}{\operatorname{\Sigma}} ||v_{x}(x,y),v_{y}(x,y)||_{2}</math><br />
<br />
Here, P is the pixel level set of selected interest points. They locate the local minima and maxima landmarks $\{t\}$ of $\epsilon = \{e(I_1),\dots,e(I_t)\}$ and for each two consecutive landmarks create a video fragment $s$ by extracting the frames $s = \{I_{t-1},\dots,I_t\}$.<br />
<br />
[[File:golfswing.png]]<br />
<br />
===Alternative Layer===<br />
This is a key layer consisting of a standard volumetric convolutional layer followed by a designed recurrent layer. Volumetric convolutional extracts features from local neighborhoods and a recurrent layer is applied to the output and it proceeds iteratively for T times. The input of a unit at position (x,y,z) in the jth feature map of the ith AL in time t is given by,<br />
<br />
:<math>u_{ij}^{xyz}(t) = u_{ij}^{xyz}(0) + f(w_{ij}^{r}u_{ij}^{xyz}(t-1)) + b_{ij} \\ <br />
u_{ij}^{xyz}(0) = f(w_{i-1}^{c}u_{(i-1)j}^{xyz}) <br />
</math><br />
<br />
U(0): feed forward output of volumetric convolutional layer. <br />
U(t-1) : recurrent input of previous time <br />
w_{k}^{c} and w_{k}^{r}: vectorized feed-forward kernels and recurrent kernels respectively <br />
f: ReLU function<br />
<br />
Figure 3 depicts this structure:<br />
[[File:unfolded.PNG|1000px]]<br />
<br />
===Volumetric Pyramid Pooling Layer===<br />
<br />
[[File:Volumetric Pyramid Pooling Layer.png|thumb|550px|Figure 2: Volumetric Pyramid Pooling Layer]]<br />
The authors have replaced the last pooling layer with a volumetric pyramid pooling layer as we need fixed-length vectors for the fully connected layers and the AL accepts video clips of arbitrary sizes and produces outputs of variable sizes. The authors have used the max pooling to pool the responses of each kernel in each volumetric bin. The outputs are kM dimensional vectors where:<br />
<br />
M: number of bins <br />
<br />
K: Number of kernels in the last alternative layer.<br />
<br />
This layer structure allows not only for arbitrary-length videos, but also arbitrary aspect ratios and scales.<br />
<br />
==Overall Architecture== <br />
[[File:DANN Architecture.png|thumb|550px|Figure 3:DANN Architecture]]<br />
The following are the components of the DANN (as shown in Figure 3)<br />
* 6 Alternative layers with 64, 128, 256, 256, 512 and 512 kernel response maps <br />
* 5 ReLU and volumetric pooling layers <br />
* 1 volumetric pyramid pooling layer <br />
* 3 fully connected layers of size 2048 each <br />
* A softmax layer<br />
<br />
==Implementation details==<br />
The authors have used the Torch toolbox platform for Implementations of volumetric convolutions, recurrent layers and optimizations. They have used a technique called as random clipping for data augmentation, in which they select a point randomly from the input video of fixed size 80x80xt after determining the temporal size t. For training the network the authors have used SGD applied to mini-batches of size 30 with a negative log likelihood criterion. Training is done by minimizing the cross-entropy loss function using backpropagation through time algorithm (BPTT). During testing, they applied a video clip divided into 80x80xt clips with a stride of 4 frames followed by testing with 10 crops. Final score is the average of all clip-level scores and the crop scores.<br />
Data augmentation techniques such as the multi-scale cropping method have been evaluated due to the recent success in the state-of-the-art performance displayed by Very Deep Two-stream ConvNets. Going by intuition, the corner cropping strategy could provide better results ( based on trade-off degree) since the receptive fields can focus harder on the central regions of the video frames [7].<br />
<br />
==Evaluations==<br />
===Datasets:===<br />
* The datasets used in the evaluation are UCF101 and HMDB51 <br />
* UCF101 – 13K videos annotated into 101 classes <br />
* HMDB51 – 6.8K videos with 51 actions. <br />
* Three training and test splits are provided <br />
* Performance measured by mean classification accuracy across the splits. <br />
* UCF101 split – 9.5K videos; HMDB51 – 3.7K training videos.<br />
<br />
===Quantitative Results===<br />
The authors used three types of optical flows, viz., sparse, RGB and TVL1 and found that TVL1 is suitable as action recognition is more easy to learn from motion information compared to raw pixel values. The influence of data augmentation is also studied. The baseline being sliding window with 75% overlap, the authors observed that the random clipping and multi-scale clipping outperformed the baseline on the UCF101 split 1 dataset. The authors were able to prove that the adaptive temporal length was able to give a boost of 4.2% when compared with architectures that had fixed-size temporal length. Experiments were also conducted to see if the learnings done in one dataset could improve the accuracy of another dataset. Fine tuning HMDB51 from UCF101 boosted the performance from 56.4% to 62.5%. The authors also observed that increasing the AL layers improves the performance as larger contexts are being embedded into the DANN. The DANN achieved an overall accuracy of 65.9% and 91.6% on HMDB51 and UCF101 respectively.<br />
<br />
<br />
[[File:Performance Comparison of different input modalities.png]]<br />
<br />
===Qualitative Analysis===<br />
The authors have discussed the quality of the prediction in the video clips taking examples of two different scenes involving bowling and haircut. In the bowling scene, the adaptive temporal choice used by DANN could aggregate more reasonable semantic structures and hence it leveraged reasonable video clips as input. On the other hand, the performance on the haircut video clip was not up to the mark as the rich contexts provided by the DANN was not helpful in a setting with simple actions performed in a simple background.<br />
<br />
==Conclusion==<br />
* Deep alternative neural network is introduced for action recognition. <br />
* DANN consists of volumetric convolutional layer and a recurrent layer. <br />
* The authors have experimented with different datasets like HMDB51 and UCF101 with different scenarios and compared the * * performance of DANN with other approaches. <br />
* The spatial size is still chosen in an ad hoc manner and this can be an area of improvement. <br />
* There are prospects for studying action tube which is a more compact input.<br />
<br />
Github code: https://github.com/wangjinzhuo/DANN<br />
<br />
==References==<br />
<br />
[1] Andrej Karpathy, George Toderici, Sachin Shetty, Tommy Leung, Rahul Sukthankar, and Li FeiFei. Large-scale video classification with convolutional neural networks. In CVPR, pages 1725–1732, 2014 <br />
<br />
[2] Karen Simonyan and Andrew Zisserman. Two-stream convolutional networks for action recognition in videos. In NIPS, pages 568–576, 2014. <br />
<br />
[3]Limin Wang, Yu Qiao, and Xiaoou Tang. Action recognition with trajectory-pooled deepconvolutional descriptors. In CVPR, pages 4305–4314, 2015. <br />
<br />
[4] Shuiwang Ji, Wei Xu, Ming Yang, and Kai Yu. 3d convolutional neural networks for human action recognition. TPAMI, 35(1):221–231, 2013. <br />
<br />
[5] Du Tran, Lubomir Bourdev, Rob Fergus, Lorenzo Torresani, and Manohar Paluri. Learning spatiotemporal features with 3d convolutional networks. In ICCV, pages 4489–4497, 2015. <br />
<br />
[6]Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action recognition. arXiv preprint arXiv:1604.04494, 2016. <br />
<br />
[7]Limin Wang, Yuanjun Xiong, Zhe Wang, Yu Qiao. Towards Good Practices for Very Deep Two-Stream ConvNets. arXiv preprint arXiv:1507.02159 , 2015. <br />
<br />
[8] IEEE International Symposium on Multimedia 2013 <br />
<br />
[36] Christopher Zach, Thomas Pock, and Horst Bischof. A duality based approach for realtime tv-l<br />
1 optical flow. In Pattern Recognition, pages 214–223. 2007.<br />
<br />
A list of expert reviews: http://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/480.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Alternative_Neural_Network:_Exploring_Contexts_As_Early_As_Possible_For_Action_Recognition&diff=30809Deep Alternative Neural Network: Exploring Contexts As Early As Possible For Action Recognition2017-11-19T18:22:19Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>==Introduction==<br />
<br />
Action recognition deals with recognizing and classifying the actions or activities done by humans or other agents in a video clip. In action recognition, Contexts contribute semantic clues for action recognition in video(See Fig below[8]). Conventional Neural Networks [1,2,3] and their shifted version 3D CNNs [4,5,6] have been employed in action recognition but they identify and aggregate the contexts at later stages. <br />
[[File:ActionRecognition1.jpg|center|400px|border|context and action region]]<br />
<br />
The authors have come up with a strategy to identify contexts in the videos as early as possible and leverage their evolutions for action recognition. Contexts contribute semantic clues for action recognition in videos. The network themselves involve a lot of layers and the first layer typically being receptive fields (RF) output only extra local features. As we go deeper into the layers the Receptive Fields expand and we start getting the contexts. The authors identified that increasing the number of layers will only cause additional burden in terms of handling the parameters and contexts could be obtained even in the earlier stages. The authors also cite the papers [9,10] that relate the CNNs and the visual systems of our brain, one remarkable difference being the abundant recurrent connections in our brain compared to the forward connections in the CNNs. <br />
<br />
The main contributions in the paper can be summarized as follows: <br />
* A Deep Alternative Neural Network (DANN) is proposed for action recognition. <br />
* DANN consists of alternative volumetric convolutional and recurrent layers. <br />
* An adaptive method to determine the temporal size of the video clip <br />
* A volumetric pyramid pooling layer to resize the output before fully connected layers.<br />
<br />
==Deep Alternative Neural Network:==<br />
===Adaptive Network Input===<br />
The input size of the video clip is generally determined empirically and various approaches have been taken in the past with a different number of frames. For instance, many previous papers suggested to used shorter intervals of between 1 to 16 frames. However, more recent work recognized that human-based actions often “span tens or hundreds of frames” and longer intervals such as 60 frames will outperform the one with shorter interval. However, there’s still no systematic way of determining number of frames for input size of the network. This serves the motives for the authors of this paper to develop this adaptive method. Past research shows that motion energy intensity induced by human activity exhibits has a regular periodicity. This signal can be approximately estimated by optical flow computation as shown in Figure 1, and is particularly suitable to address our temporal estimation due to: <br />
* the local minima and maxima landmarks probably correspond to characteristic gesture and motion <br />
* it is relatively robust to changes in camera viewpoint.<br />
<br />
The authors have come up with an adaptive method to automatically select the most discriminative video fragments using the density of optical flow energy which exhibits regular periodicity. According to Wikipedia, optical flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene, and optical flow methods try to calculate the motion between two image frames which are taken at different times. The optimal flow energy of an optical field $(v_{x},v_{y})$ is defined as follows <br />
<br />
:<math>e(I)=\underset{(x,y)\in\mathbb{P}}{\operatorname{\Sigma}} ||v_{x}(x,y),v_{y}(x,y)||_{2}</math><br />
<br />
Here, P is the pixel level set of selected interest points. They locate the local minima and maxima landmarks $\{t\}$ of $\epsilon = \{e(I_1),\dots,e(I_t)\}$ and for each two consecutive landmarks create a video fragment $s$ by extracting the frames $s = \{I_{t-1},\dots,I_t\}$.<br />
<br />
[[File:golfswing.png]]<br />
<br />
===Alternative Layer===<br />
This is a key layer consisting of a standard volumetric convolutional layer followed by a designed recurrent layer. Volumetric convolutional extracts features from local neighborhoods and a recurrent layer is applied to the output and it proceeds iteratively for T times. The input of a unit at position (x,y,z) in the jth feature map of the ith AL in time t is given by,<br />
<br />
:<math>u_{ij}^{xyz}(t) = u_{ij}^{xyz}(0) + f(w_{ij}^{r}u_{ij}^{xyz}(t-1)) + b_{ij} \\ <br />
u_{ij}^{xyz}(0) = f(w_{i-1}^{c}u_{(i-1)j}^{xyz}) <br />
</math><br />
<br />
U(0): feed forward output of volumetric convolutional layer. <br />
U(t-1) : recurrent input of previous time <br />
w_{k}^{c} and w_{k}^{r}: vectorized feed-forward kernels and recurrent kernels respectively <br />
f: ReLU function<br />
<br />
Figure 3 depicts this structure:<br />
[[File:unfolded.PNG|1000px]]<br />
<br />
===Volumetric Pyramid Pooling Layer===<br />
<br />
[[File:Volumetric Pyramid Pooling Layer.png|thumb|550px|Figure 2: Volumetric Pyramid Pooling Layer]]<br />
The authors have replaced the last pooling layer with a volumetric pyramid pooling layer as we need fixed-length vectors for the fully connected layers and the AL accepts video clips of arbitrary sizes and produces outputs of variable sizes. The authors have used the max pooling to pool the responses of each kernel in each volumetric bin. The outputs are kM dimensional vectors where:<br />
<br />
M: number of bins <br />
<br />
K: Number of kernels in the last alternative layer.<br />
<br />
This layer structure allows not only for arbitrary-length videos, but also arbitrary aspect ratios and scales.<br />
<br />
==Overall Architecture== <br />
[[File:DANN Architecture.png|thumb|550px|Figure 3:DANN Architecture]]<br />
The following are the components of the DANN (as shown in Figure 3)<br />
* 6 Alternative layers with 64, 128, 256, 256, 512 and 512 kernel response maps <br />
* 5 ReLU and volumetric pooling layers <br />
* 1 volumetric pyramid pooling layer <br />
* 3 fully connected layers of size 2048 each <br />
* A softmax layer<br />
<br />
==Implementation details==<br />
The authors have used the Torch toolbox platform for Implementations of volumetric convolutions, recurrent layers and optimizations. They have used a technique called as random clipping for data augmentation, in which they select a point randomly from the input video of fixed size 80x80xt after determining the temporal size t. For training the network the authors have used SGD applied to mini-batches of size 30 with a negative log likelihood criterion. Training is done by minimizing the cross-entropy loss function using backpropagation through time algorithm (BPTT). During testing, they applied a video clip divided into 80x80xt clips with a stride of 4 frames followed by testing with 10 crops. Final score is the average of all clip-level scores and the crop scores.<br />
Data augmentation techniques such as the multi-scale cropping method have been evaluated due to the recent success in the state-of-the-art performance displayed by Very Deep Two-stream ConvNets. Going by intuition, the corner cropping strategy could provide better results ( based on trade-off degree) since the receptive fields can focus harder on the central regions of the video frames [7].<br />
<br />
==Evaluations==<br />
===Datasets:===<br />
* The datasets used in the evaluation are UCF101 and HMDB51 <br />
* UCF101 – 13K videos annotated into 101 classes <br />
* HMDB51 – 6.8K videos with 51 actions. <br />
* Three training and test splits are provided <br />
* Performance measured by mean classification accuracy across the splits. <br />
* UCF101 split – 9.5K videos; HMDB51 – 3.7K training videos.<br />
<br />
===Quantitative Results===<br />
The authors used three types of optical flows, viz., sparse, RGB and TVL1 and found that TVL1 is suitable as action recognition is more easy to learn from motion information compared to raw pixel values. The influence of data augmentation is also studied. The baseline being sliding window with 75% overlap, the authors observed that the random clipping and multi-scale clipping outperformed the baseline on the UCF101 split 1 dataset. The authors were able to prove that the adaptive temporal length was able to give a boost of 4.2% when compared with architectures that had fixed-size temporal length. Experiments were also conducted to see if the learnings done in one dataset could improve the accuracy of another dataset. Fine tuning HMDB51 from UCF101 boosted the performance from 56.4% to 62.5%. The authors also observed that increasing the AL layers improves the performance as larger contexts are being embedded into the DANN. The DANN achieved an overall accuracy of 65.9% and 91.6% on HMDB51 and UCF101 respectively.<br />
<br />
<br />
[[File:Performance Comparison of different input modalities.png]]<br />
<br />
===Qualitative Analysis===<br />
The authors have discussed the quality of the prediction in the video clips taking examples of two different scenes involving bowling and haircut. In the bowling scene, the adaptive temporal choice used by DANN could aggregate more reasonable semantic structures and hence it leveraged reasonable video clips as input. On the other hand, the performance on the haircut video clip was not up to the mark as the rich contexts provided by the DANN was not helpful in a setting with simple actions performed in a simple background.<br />
<br />
==Conclusion==<br />
* Deep alternative neural network is introduced for action recognition. <br />
* DANN consists of volumetric convolutional layer and a recurrent layer. <br />
* The authors have experimented with different datasets like HMDB51 and UCF101 with different scenarios and compared the * * performance of DANN with other approaches. <br />
* The spatial size is still chosen in an ad hoc manner and this can be an area of improvement. <br />
* There are prospects for studying action tube which is a more compact input.<br />
<br />
Github code: https://github.com/wangjinzhuo/DANN<br />
<br />
==References==<br />
<br />
[1] Andrej Karpathy, George Toderici, Sachin Shetty, Tommy Leung, Rahul Sukthankar, and Li FeiFei. Large-scale video classification with convolutional neural networks. In CVPR, pages 1725–1732, 2014 <br />
<br />
[2] Karen Simonyan and Andrew Zisserman. Two-stream convolutional networks for action recognition in videos. In NIPS, pages 568–576, 2014. <br />
<br />
[3]Limin Wang, Yu Qiao, and Xiaoou Tang. Action recognition with trajectory-pooled deepconvolutional descriptors. In CVPR, pages 4305–4314, 2015. <br />
<br />
[4] Shuiwang Ji, Wei Xu, Ming Yang, and Kai Yu. 3d convolutional neural networks for human action recognition. TPAMI, 35(1):221–231, 2013. <br />
<br />
[5] Du Tran, Lubomir Bourdev, Rob Fergus, Lorenzo Torresani, and Manohar Paluri. Learning spatiotemporal features with 3d convolutional networks. In ICCV, pages 4489–4497, 2015. <br />
<br />
[6]Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action recognition. arXiv preprint arXiv:1604.04494, 2016. <br />
<br />
[7]Limin Wang, Yuanjun Xiong, Zhe Wang, Yu Qiao. Towards Good Practices for Very Deep Two-Stream ConvNets. arXiv preprint arXiv:1507.02159 , 2015. <br />
<br />
[36] Christopher Zach, Thomas Pock, and Horst Bischof. A duality based approach for realtime tv-l<br />
1 optical flow. In Pattern Recognition, pages 214–223. 2007.<br />
<br />
A list of expert reviews: http://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/480.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Alternative_Neural_Network:_Exploring_Contexts_As_Early_As_Possible_For_Action_Recognition&diff=30808Deep Alternative Neural Network: Exploring Contexts As Early As Possible For Action Recognition2017-11-19T18:22:01Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>==Introduction==<br />
<br />
Action recognition deals with recognizing and classifying the actions or activities done by humans or other agents in a video clip. In action recognition, Contexts contribute semantic clues for action recognition in video(See Fig below). Conventional Neural Networks [1,2,3] and their shifted version 3D CNNs [4,5,6] have been employed in action recognition but they identify and aggregate the contexts at later stages. <br />
[[File:ActionRecognition1.jpg|center|400px|border|context and action region]]<br />
<br />
The authors have come up with a strategy to identify contexts in the videos as early as possible and leverage their evolutions for action recognition. Contexts contribute semantic clues for action recognition in videos. The network themselves involve a lot of layers and the first layer typically being receptive fields (RF) output only extra local features. As we go deeper into the layers the Receptive Fields expand and we start getting the contexts. The authors identified that increasing the number of layers will only cause additional burden in terms of handling the parameters and contexts could be obtained even in the earlier stages. The authors also cite the papers [9,10] that relate the CNNs and the visual systems of our brain, one remarkable difference being the abundant recurrent connections in our brain compared to the forward connections in the CNNs. <br />
<br />
The main contributions in the paper can be summarized as follows: <br />
* A Deep Alternative Neural Network (DANN) is proposed for action recognition. <br />
* DANN consists of alternative volumetric convolutional and recurrent layers. <br />
* An adaptive method to determine the temporal size of the video clip <br />
* A volumetric pyramid pooling layer to resize the output before fully connected layers.<br />
<br />
==Deep Alternative Neural Network:==<br />
===Adaptive Network Input===<br />
The input size of the video clip is generally determined empirically and various approaches have been taken in the past with a different number of frames. For instance, many previous papers suggested to used shorter intervals of between 1 to 16 frames. However, more recent work recognized that human-based actions often “span tens or hundreds of frames” and longer intervals such as 60 frames will outperform the one with shorter interval. However, there’s still no systematic way of determining number of frames for input size of the network. This serves the motives for the authors of this paper to develop this adaptive method. Past research shows that motion energy intensity induced by human activity exhibits has a regular periodicity. This signal can be approximately estimated by optical flow computation as shown in Figure 1, and is particularly suitable to address our temporal estimation due to: <br />
* the local minima and maxima landmarks probably correspond to characteristic gesture and motion <br />
* it is relatively robust to changes in camera viewpoint.<br />
<br />
The authors have come up with an adaptive method to automatically select the most discriminative video fragments using the density of optical flow energy which exhibits regular periodicity. According to Wikipedia, optical flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene, and optical flow methods try to calculate the motion between two image frames which are taken at different times. The optimal flow energy of an optical field $(v_{x},v_{y})$ is defined as follows <br />
<br />
:<math>e(I)=\underset{(x,y)\in\mathbb{P}}{\operatorname{\Sigma}} ||v_{x}(x,y),v_{y}(x,y)||_{2}</math><br />
<br />
Here, P is the pixel level set of selected interest points. They locate the local minima and maxima landmarks $\{t\}$ of $\epsilon = \{e(I_1),\dots,e(I_t)\}$ and for each two consecutive landmarks create a video fragment $s$ by extracting the frames $s = \{I_{t-1},\dots,I_t\}$.<br />
<br />
[[File:golfswing.png]]<br />
<br />
===Alternative Layer===<br />
This is a key layer consisting of a standard volumetric convolutional layer followed by a designed recurrent layer. Volumetric convolutional extracts features from local neighborhoods and a recurrent layer is applied to the output and it proceeds iteratively for T times. The input of a unit at position (x,y,z) in the jth feature map of the ith AL in time t is given by,<br />
<br />
:<math>u_{ij}^{xyz}(t) = u_{ij}^{xyz}(0) + f(w_{ij}^{r}u_{ij}^{xyz}(t-1)) + b_{ij} \\ <br />
u_{ij}^{xyz}(0) = f(w_{i-1}^{c}u_{(i-1)j}^{xyz}) <br />
</math><br />
<br />
U(0): feed forward output of volumetric convolutional layer. <br />
U(t-1) : recurrent input of previous time <br />
w_{k}^{c} and w_{k}^{r}: vectorized feed-forward kernels and recurrent kernels respectively <br />
f: ReLU function<br />
<br />
Figure 3 depicts this structure:<br />
[[File:unfolded.PNG|1000px]]<br />
<br />
===Volumetric Pyramid Pooling Layer===<br />
<br />
[[File:Volumetric Pyramid Pooling Layer.png|thumb|550px|Figure 2: Volumetric Pyramid Pooling Layer]]<br />
The authors have replaced the last pooling layer with a volumetric pyramid pooling layer as we need fixed-length vectors for the fully connected layers and the AL accepts video clips of arbitrary sizes and produces outputs of variable sizes. The authors have used the max pooling to pool the responses of each kernel in each volumetric bin. The outputs are kM dimensional vectors where:<br />
<br />
M: number of bins <br />
<br />
K: Number of kernels in the last alternative layer.<br />
<br />
This layer structure allows not only for arbitrary-length videos, but also arbitrary aspect ratios and scales.<br />
<br />
==Overall Architecture== <br />
[[File:DANN Architecture.png|thumb|550px|Figure 3:DANN Architecture]]<br />
The following are the components of the DANN (as shown in Figure 3)<br />
* 6 Alternative layers with 64, 128, 256, 256, 512 and 512 kernel response maps <br />
* 5 ReLU and volumetric pooling layers <br />
* 1 volumetric pyramid pooling layer <br />
* 3 fully connected layers of size 2048 each <br />
* A softmax layer<br />
<br />
==Implementation details==<br />
The authors have used the Torch toolbox platform for Implementations of volumetric convolutions, recurrent layers and optimizations. They have used a technique called as random clipping for data augmentation, in which they select a point randomly from the input video of fixed size 80x80xt after determining the temporal size t. For training the network the authors have used SGD applied to mini-batches of size 30 with a negative log likelihood criterion. Training is done by minimizing the cross-entropy loss function using backpropagation through time algorithm (BPTT). During testing, they applied a video clip divided into 80x80xt clips with a stride of 4 frames followed by testing with 10 crops. Final score is the average of all clip-level scores and the crop scores.<br />
Data augmentation techniques such as the multi-scale cropping method have been evaluated due to the recent success in the state-of-the-art performance displayed by Very Deep Two-stream ConvNets. Going by intuition, the corner cropping strategy could provide better results ( based on trade-off degree) since the receptive fields can focus harder on the central regions of the video frames [7].<br />
<br />
==Evaluations==<br />
===Datasets:===<br />
* The datasets used in the evaluation are UCF101 and HMDB51 <br />
* UCF101 – 13K videos annotated into 101 classes <br />
* HMDB51 – 6.8K videos with 51 actions. <br />
* Three training and test splits are provided <br />
* Performance measured by mean classification accuracy across the splits. <br />
* UCF101 split – 9.5K videos; HMDB51 – 3.7K training videos.<br />
<br />
===Quantitative Results===<br />
The authors used three types of optical flows, viz., sparse, RGB and TVL1 and found that TVL1 is suitable as action recognition is more easy to learn from motion information compared to raw pixel values. The influence of data augmentation is also studied. The baseline being sliding window with 75% overlap, the authors observed that the random clipping and multi-scale clipping outperformed the baseline on the UCF101 split 1 dataset. The authors were able to prove that the adaptive temporal length was able to give a boost of 4.2% when compared with architectures that had fixed-size temporal length. Experiments were also conducted to see if the learnings done in one dataset could improve the accuracy of another dataset. Fine tuning HMDB51 from UCF101 boosted the performance from 56.4% to 62.5%. The authors also observed that increasing the AL layers improves the performance as larger contexts are being embedded into the DANN. The DANN achieved an overall accuracy of 65.9% and 91.6% on HMDB51 and UCF101 respectively.<br />
<br />
<br />
[[File:Performance Comparison of different input modalities.png]]<br />
<br />
===Qualitative Analysis===<br />
The authors have discussed the quality of the prediction in the video clips taking examples of two different scenes involving bowling and haircut. In the bowling scene, the adaptive temporal choice used by DANN could aggregate more reasonable semantic structures and hence it leveraged reasonable video clips as input. On the other hand, the performance on the haircut video clip was not up to the mark as the rich contexts provided by the DANN was not helpful in a setting with simple actions performed in a simple background.<br />
<br />
==Conclusion==<br />
* Deep alternative neural network is introduced for action recognition. <br />
* DANN consists of volumetric convolutional layer and a recurrent layer. <br />
* The authors have experimented with different datasets like HMDB51 and UCF101 with different scenarios and compared the * * performance of DANN with other approaches. <br />
* The spatial size is still chosen in an ad hoc manner and this can be an area of improvement. <br />
* There are prospects for studying action tube which is a more compact input.<br />
<br />
Github code: https://github.com/wangjinzhuo/DANN<br />
<br />
==References==<br />
<br />
[1] Andrej Karpathy, George Toderici, Sachin Shetty, Tommy Leung, Rahul Sukthankar, and Li FeiFei. Large-scale video classification with convolutional neural networks. In CVPR, pages 1725–1732, 2014 <br />
<br />
[2] Karen Simonyan and Andrew Zisserman. Two-stream convolutional networks for action recognition in videos. In NIPS, pages 568–576, 2014. <br />
<br />
[3]Limin Wang, Yu Qiao, and Xiaoou Tang. Action recognition with trajectory-pooled deepconvolutional descriptors. In CVPR, pages 4305–4314, 2015. <br />
<br />
[4] Shuiwang Ji, Wei Xu, Ming Yang, and Kai Yu. 3d convolutional neural networks for human action recognition. TPAMI, 35(1):221–231, 2013. <br />
<br />
[5] Du Tran, Lubomir Bourdev, Rob Fergus, Lorenzo Torresani, and Manohar Paluri. Learning spatiotemporal features with 3d convolutional networks. In ICCV, pages 4489–4497, 2015. <br />
<br />
[6]Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action recognition. arXiv preprint arXiv:1604.04494, 2016. <br />
<br />
[7]Limin Wang, Yuanjun Xiong, Zhe Wang, Yu Qiao. Towards Good Practices for Very Deep Two-Stream ConvNets. arXiv preprint arXiv:1507.02159 , 2015. <br />
<br />
[36] Christopher Zach, Thomas Pock, and Horst Bischof. A duality based approach for realtime tv-l<br />
1 optical flow. In Pattern Recognition, pages 214–223. 2007.<br />
<br />
A list of expert reviews: http://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/480.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=File:ActionRecognition1.jpg&diff=30807File:ActionRecognition1.jpg2017-11-19T18:19:07Z<p>Nkadiyar: Action Recognition in Videos</p>
<hr />
<div>Action Recognition in Videos</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Deep_Alternative_Neural_Network:_Exploring_Contexts_As_Early_As_Possible_For_Action_Recognition&diff=30806Deep Alternative Neural Network: Exploring Contexts As Early As Possible For Action Recognition2017-11-19T18:15:25Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>==Introduction==<br />
<br />
Action recognition deals with recognizing and classifying the actions or activities done by humans or other agents in a video clip. In action recognition, Contexts contribute semantic clues for action recognition in video. Conventional Neural Networks [1,2,3] and their shifted version 3D CNNs [4,5,6] have been employed in action recognition but they identify and aggregate the contexts at later stages. <br />
The authors have come up with a strategy to identify contexts in the videos as early as possible and leverage their evolutions for action recognition. Contexts contribute semantic clues for action recognition in videos. The network themselves involve a lot of layers and the first layer typically being receptive fields (RF) output only extra local features. As we go deeper into the layers the Receptive Fields expand and we start getting the contexts. The authors identified that increasing the number of layers will only cause additional burden in terms of handling the parameters and contexts could be obtained even in the earlier stages. The authors also cite the papers [9,10] that relate the CNNs and the visual systems of our brain, one remarkable difference being the abundant recurrent connections in our brain compared to the forward connections in the CNNs. <br />
<br />
The main contributions in the paper can be summarized as follows: <br />
* A Deep Alternative Neural Network (DANN) is proposed for action recognition. <br />
* DANN consists of alternative volumetric convolutional and recurrent layers. <br />
* An adaptive method to determine the temporal size of the video clip <br />
* A volumetric pyramid pooling layer to resize the output before fully connected layers.<br />
<br />
==Deep Alternative Neural Network:==<br />
===Adaptive Network Input===<br />
The input size of the video clip is generally determined empirically and various approaches have been taken in the past with a different number of frames. For instance, many previous papers suggested to used shorter intervals of between 1 to 16 frames. However, more recent work recognized that human-based actions often “span tens or hundreds of frames” and longer intervals such as 60 frames will outperform the one with shorter interval. However, there’s still no systematic way of determining number of frames for input size of the network. This serves the motives for the authors of this paper to develop this adaptive method. Past research shows that motion energy intensity induced by human activity exhibits has a regular periodicity. This signal can be approximately estimated by optical flow computation as shown in Figure 1, and is particularly suitable to address our temporal estimation due to: <br />
* the local minima and maxima landmarks probably correspond to characteristic gesture and motion <br />
* it is relatively robust to changes in camera viewpoint.<br />
<br />
The authors have come up with an adaptive method to automatically select the most discriminative video fragments using the density of optical flow energy which exhibits regular periodicity. According to Wikipedia, optical flow is the pattern of apparent motion of objects, surfaces, and edges in a visual scene caused by the relative motion between an observer and a scene, and optical flow methods try to calculate the motion between two image frames which are taken at different times. The optimal flow energy of an optical field $(v_{x},v_{y})$ is defined as follows <br />
<br />
:<math>e(I)=\underset{(x,y)\in\mathbb{P}}{\operatorname{\Sigma}} ||v_{x}(x,y),v_{y}(x,y)||_{2}</math><br />
<br />
Here, P is the pixel level set of selected interest points. They locate the local minima and maxima landmarks $\{t\}$ of $\epsilon = \{e(I_1),\dots,e(I_t)\}$ and for each two consecutive landmarks create a video fragment $s$ by extracting the frames $s = \{I_{t-1},\dots,I_t\}$.<br />
<br />
[[File:golfswing.png]]<br />
<br />
===Alternative Layer===<br />
This is a key layer consisting of a standard volumetric convolutional layer followed by a designed recurrent layer. Volumetric convolutional extracts features from local neighborhoods and a recurrent layer is applied to the output and it proceeds iteratively for T times. The input of a unit at position (x,y,z) in the jth feature map of the ith AL in time t is given by,<br />
<br />
:<math>u_{ij}^{xyz}(t) = u_{ij}^{xyz}(0) + f(w_{ij}^{r}u_{ij}^{xyz}(t-1)) + b_{ij} \\ <br />
u_{ij}^{xyz}(0) = f(w_{i-1}^{c}u_{(i-1)j}^{xyz}) <br />
</math><br />
<br />
U(0): feed forward output of volumetric convolutional layer. <br />
U(t-1) : recurrent input of previous time <br />
w_{k}^{c} and w_{k}^{r}: vectorized feed-forward kernels and recurrent kernels respectively <br />
f: ReLU function<br />
<br />
Figure 3 depicts this structure:<br />
[[File:unfolded.PNG|1000px]]<br />
<br />
===Volumetric Pyramid Pooling Layer===<br />
<br />
[[File:Volumetric Pyramid Pooling Layer.png|thumb|550px|Figure 2: Volumetric Pyramid Pooling Layer]]<br />
The authors have replaced the last pooling layer with a volumetric pyramid pooling layer as we need fixed-length vectors for the fully connected layers and the AL accepts video clips of arbitrary sizes and produces outputs of variable sizes. The authors have used the max pooling to pool the responses of each kernel in each volumetric bin. The outputs are kM dimensional vectors where:<br />
<br />
M: number of bins <br />
<br />
K: Number of kernels in the last alternative layer.<br />
<br />
This layer structure allows not only for arbitrary-length videos, but also arbitrary aspect ratios and scales.<br />
<br />
==Overall Architecture== <br />
[[File:DANN Architecture.png|thumb|550px|Figure 3:DANN Architecture]]<br />
The following are the components of the DANN (as shown in Figure 3)<br />
* 6 Alternative layers with 64, 128, 256, 256, 512 and 512 kernel response maps <br />
* 5 ReLU and volumetric pooling layers <br />
* 1 volumetric pyramid pooling layer <br />
* 3 fully connected layers of size 2048 each <br />
* A softmax layer<br />
<br />
==Implementation details==<br />
The authors have used the Torch toolbox platform for Implementations of volumetric convolutions, recurrent layers and optimizations. They have used a technique called as random clipping for data augmentation, in which they select a point randomly from the input video of fixed size 80x80xt after determining the temporal size t. For training the network the authors have used SGD applied to mini-batches of size 30 with a negative log likelihood criterion. Training is done by minimizing the cross-entropy loss function using backpropagation through time algorithm (BPTT). During testing, they applied a video clip divided into 80x80xt clips with a stride of 4 frames followed by testing with 10 crops. Final score is the average of all clip-level scores and the crop scores.<br />
Data augmentation techniques such as the multi-scale cropping method have been evaluated due to the recent success in the state-of-the-art performance displayed by Very Deep Two-stream ConvNets. Going by intuition, the corner cropping strategy could provide better results ( based on trade-off degree) since the receptive fields can focus harder on the central regions of the video frames [7].<br />
<br />
==Evaluations==<br />
===Datasets:===<br />
* The datasets used in the evaluation are UCF101 and HMDB51 <br />
* UCF101 – 13K videos annotated into 101 classes <br />
* HMDB51 – 6.8K videos with 51 actions. <br />
* Three training and test splits are provided <br />
* Performance measured by mean classification accuracy across the splits. <br />
* UCF101 split – 9.5K videos; HMDB51 – 3.7K training videos.<br />
<br />
===Quantitative Results===<br />
The authors used three types of optical flows, viz., sparse, RGB and TVL1 and found that TVL1 is suitable as action recognition is more easy to learn from motion information compared to raw pixel values. The influence of data augmentation is also studied. The baseline being sliding window with 75% overlap, the authors observed that the random clipping and multi-scale clipping outperformed the baseline on the UCF101 split 1 dataset. The authors were able to prove that the adaptive temporal length was able to give a boost of 4.2% when compared with architectures that had fixed-size temporal length. Experiments were also conducted to see if the learnings done in one dataset could improve the accuracy of another dataset. Fine tuning HMDB51 from UCF101 boosted the performance from 56.4% to 62.5%. The authors also observed that increasing the AL layers improves the performance as larger contexts are being embedded into the DANN. The DANN achieved an overall accuracy of 65.9% and 91.6% on HMDB51 and UCF101 respectively.<br />
<br />
<br />
[[File:Performance Comparison of different input modalities.png]]<br />
<br />
===Qualitative Analysis===<br />
The authors have discussed the quality of the prediction in the video clips taking examples of two different scenes involving bowling and haircut. In the bowling scene, the adaptive temporal choice used by DANN could aggregate more reasonable semantic structures and hence it leveraged reasonable video clips as input. On the other hand, the performance on the haircut video clip was not up to the mark as the rich contexts provided by the DANN was not helpful in a setting with simple actions performed in a simple background.<br />
<br />
==Conclusion==<br />
* Deep alternative neural network is introduced for action recognition. <br />
* DANN consists of volumetric convolutional layer and a recurrent layer. <br />
* The authors have experimented with different datasets like HMDB51 and UCF101 with different scenarios and compared the * * performance of DANN with other approaches. <br />
* The spatial size is still chosen in an ad hoc manner and this can be an area of improvement. <br />
* There are prospects for studying action tube which is a more compact input.<br />
<br />
Github code: https://github.com/wangjinzhuo/DANN<br />
<br />
==References==<br />
<br />
[1] Andrej Karpathy, George Toderici, Sachin Shetty, Tommy Leung, Rahul Sukthankar, and Li FeiFei. Large-scale video classification with convolutional neural networks. In CVPR, pages 1725–1732, 2014 <br />
<br />
[2] Karen Simonyan and Andrew Zisserman. Two-stream convolutional networks for action recognition in videos. In NIPS, pages 568–576, 2014. <br />
<br />
[3]Limin Wang, Yu Qiao, and Xiaoou Tang. Action recognition with trajectory-pooled deepconvolutional descriptors. In CVPR, pages 4305–4314, 2015. <br />
<br />
[4] Shuiwang Ji, Wei Xu, Ming Yang, and Kai Yu. 3d convolutional neural networks for human action recognition. TPAMI, 35(1):221–231, 2013. <br />
<br />
[5] Du Tran, Lubomir Bourdev, Rob Fergus, Lorenzo Torresani, and Manohar Paluri. Learning spatiotemporal features with 3d convolutional networks. In ICCV, pages 4489–4497, 2015. <br />
<br />
[6]Gül Varol, Ivan Laptev, and Cordelia Schmid. Long-term temporal convolutions for action recognition. arXiv preprint arXiv:1604.04494, 2016. <br />
<br />
[7]Limin Wang, Yuanjun Xiong, Zhe Wang, Yu Qiao. Towards Good Practices for Very Deep Two-Stream ConvNets. arXiv preprint arXiv:1507.02159 , 2015. <br />
<br />
[36] Christopher Zach, Thomas Pock, and Horst Bischof. A duality based approach for realtime tv-l<br />
1 optical flow. In Pattern Recognition, pages 214–223. 2007.<br />
<br />
A list of expert reviews: http://media.nips.cc/nipsbooks/nipspapers/paper_files/nips29/reviews/480.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=LightRNN:_Memory_and_Computation-Efficient_Recurrent_Neural_Networks&diff=30444LightRNN: Memory and Computation-Efficient Recurrent Neural Networks2017-11-15T18:22:52Z<p>Nkadiyar: /* Part I: 2-Component Shared Embedding */</p>
<hr />
<div>= Introduction =<br />
<br />
The study of natural language processing has been around for more than fifty years. It begins in 1950s which the specific field of natural language processing (NLP) is still embedded in the subject of linguistics (Hirschberg & Manning, 2015). After the emergence of strong computational power, computational linguistics began to evolve and gradually branch out to various applications in NLP, such as text classification, speech recognition and question answering (Brownlee, 2017). Computational linguistics or natural language processing is usually defined as “subfield of computer science concerned with using computational techniques to learn, understand, and produce human language content” (Hirschberg & Manning, 2015, p. 261). <br />
<br />
With the development of deep neural networks, one type of neural network, namely recurrent neural networks (RNN) have performed significantly well in many natural language processing tasks. The reason is that nature of RNN takes into account the past inputs as well as the current input without resulting in vanishing or exploding gradient. More detail of how RNN works in the context of NLP will be discussed in the section of NLP using RNN. However, one limitation of RNN used in NLP is its enormous size of input vocabulary. This will result in a very complex RNN model with too many parameters to train and makes the training process both time and memory-consuming. This serves as the major motivation for this paper’s authors to develop a new technique utilized in RNN, which is particularly efficient at processing large size of vocabulary in many NLP tasks, namely LightRNN.<br />
<br />
= Motivations =<br />
<br />
In language modelling, researchers used to represent words by arbitrary codes, such as “Id143” is the code for “dog” (“Vector Representations of Words,” 2017). Such coding of words is completely random, and it loses the meaning of the words and (more importantly) connection with other words. Nowadays, one-hot representation of words is commonly used, in which a word is represented by a vector of numbers and the dimension of the vector is related to the size of the vocabulary. In RNN, all words in the vocabulary are coded using one-hot representation and then mapped to an embedding vector (Li, Qin, Yang, Hu, & Liu, 2016). Such embedding vector is “a continuous vector space where semantically similar words are mapped to nearby points” (“Vector Representations of Words” 2017, para. 6). Popular RNN structure used in NLP task is long short-term memory (LSTM). In order to predict the probability of the next word, the last hidden layer of the network needs to calculate the probability distribution over all other words in the vocabulary. Note that the most time-consuming operation in RNNs is to calculate a probability distribution over all the words in the vocabulary, which requires the multiplication of the output-embedding matrix and the hidden state at each position of a sequence. Lastly, an activation function (commonly, softmax function) is used to select the next word with the highest probability. <br />
This method has 3 major limitations:<br />
<br />
1. Memory Constraint <br />
<br />
When input vocabulary contains an enormous amount of unique words, which is very common in various NLP tasks, the size of model becomes very large. This means the number of trainable parameters is very big, which makes it difficult to fit such model on a regular GPU device.<br />
<br />
2. Computationally Heavy to Train<br />
<br />
As previously mentioned, the probability distribution of all other words in the vocabulary needs to be computed to determine what predicted word it would be. When the size of the vocabulary is large, such calculation can be computationally heavy. <br />
<br />
3. Low Compressability<br />
<br />
Due to the memory and computation-consuming process of RNN applied in NLP tasks, mobile devices cannot usually handle such algorithm, which makes it undesirable and limits its usage.<br />
<br />
An alternative approach to handle the overhead is by leveraging weak lower-level learners by Boosting. But a drawback is that this technique has only been implemented for a few specific tasks in the past such as Time Series predictions [Boné et. al.].<br />
<br />
= LightRNN Structure =<br />
<br />
The authors of the paper proposed a new structure that effectively reduces the size of the model by arranging all words in the vocabulary into a word table, which is referred as “2-Component (2C) shared embedding for word representation”. This is done by factorizing a vocabulary's embedding into two shared components (row and column). Thus, a word is indexed by its location in such table, which in terms is characterized by the corresponding row and column components. Each row and column component are unique row vector and column vector respectively. By organizing each word in the vocabulary in this manner, multiple words can share the same row component or column component and it can reduce the number of trainable parameters significantly. <br />
The next question is how to construct such word table. More specifically, how to allocate each word in the vocabulary to different positions so that semantically similar words are in the same row or column. The authors proposed a bootstrap method to solve this problem. Essentially, we first randomly distribute words into the table. Then, we let the model “learn” better position of each word by minimizing training error. By repeating this process, each word can be allocated to a particular position within the table so that similar words share common row or column components. More details of those 2 parts of LightRNN structure will be discussed in the following sections.<br />
<br />
There are 2 major benefits of the proposed technique:<br />
<br />
1. Computationally efficient<br />
<br />
The name “LightRNN” is to illustrate the small model size and fast training speed. Because of these features of the new RNN architecture, it’s possible to launch such model onto regular GPU and other mobile devices. <br />
<br />
2. Higher scalability <br />
<br />
The authors briefly explained this algorithm is scalable because if parallel-computing is needed to train such model, the difficulty of combining smaller models is low. <br />
<br />
<br />
== Part I: 2-Component Shared Embedding ==<br />
<br />
The key aspect of LightRNN structure is its innovative method of word representation, namely 2-Component Shared Embedding. All words in the vocabulary are organized into a table with row components and column components. Each pair of the element in a row component and a column component is corresponding to a unique word in the vocabulary. For instance, the <math>i^{th}</math> row and <math>j^{th}</math> column are the row and column indexes for <math>X_{ij}</math>. As shown in the following graph, <math>x_{1}</math> is corresponding to the words “January”. In 2C shared embedding table, it’s indexed by 2 elements: <math>x^{r}_{1}</math> and <math>x^{c}_{1}</math> where the subscript indicates which row component and column component this word belongs to. Ideally, words that share similar semantic features should be assigned to the same row or column. The shared embedding word table in Figure 1 serves as a good example: the word “one” and “January” are assigned to the same column, while the word “one” and “two” are allocated to the same row. <br />
<br />
[[File:2C shared embedding.png|700px|thumb|centre|Fig 1. 2-Component Shared Embedding for Word Representation]]<br />
<br />
The main advantage of using such word representation is it reduces the number of vector/element needed for input word embedding. For instance, if there are 25 unique words in the vocabulary, the number of vectors to represent all the words is 10, namely 5 row vectors/elements and 5 column vectors/elements. Therefore, the shared embedding word table is a 5 by 5 matrix. In general, the formula for calculating number of vector/element needed to represent <math>|V|</math> words is <math>2\sqrt{|V|}</math>.<br />
<br />
== Part II: How 2C Shared Embedding is Used in LightRNN ==<br />
<br />
After constructing such word representation table, those 2-component shared embedding matrices are fed into the recurrent neural network. The following Figure 2 demonstrates a portion of LightRNN structure (left) with comparison with the regular RNN (right). Compared to regular RNN where a single input <math>x_{t-1}</math> is fed into the network each time, 2 elements of a single input <math>x_{t-1}</math>: <math>x^{r}_{t-1}</math> and <math>x^{c}_{t-1}</math> are fed into LightRNN. <br />
<br />
[[File:LightRNN.PNG |700px|thumb|centre|Fig 2. LightRNN Structure & Regular RNN]]<br />
<br />
As mentioned before, the last hidden layer will produce the probabilities of <math>word_{t}</math>. Based on the diagram below, the following formulas are used:<br />
Let n = dimension/length of a row input vector/a column input vector, <math>X^{c}, X^{r} \in \mathbb{R}^{n \times \sqrt{|V|}}</math> denotes the input-embedding matrices: <br />
<center><br />
: row vector <math>x^{r}_{t-1} \in \mathbb{R}^n</math><br />
: column vector <math>x^{c}_{t-1} \in \mathbb{R}^n</math><br />
</center><br />
<br />
Let <math>h^{c}_{t-1}, h^{r}_{t-1} \in \mathbb{R}^m</math> denotes the two hidden layers where m = dimension of the hidden layer:<br />
<center><br />
: <math>h^{c}_{t-1} = f(W x_{t-1}^{c} + U h_{t-1}^{r} + b) </math><br />
: <math>h^{r}_{t} = f(W x_{t}^{r} + U h_{t-1}^{c} + b) </math><br />
</center><br />
where <math>W \in \mathbb{R}^{m \times n}</math>, <math>U \in \mathbb{R}^{m \times m}</math>, and <math>b \in \mathbb{R}^m</math> and <math>f</math> is a nonlinear activation function<br />
<br />
The final step in LightRNN is to calculate <math>P_{r}(w_{t})</math> and <math>P_{c}(w_{t})</math> , which means the probability of a word w at time t, using the following formulas:<br />
<center><br />
: <math>P_{r}(w_t) = \frac{exp(h_{t-1}^{c} y_{r(w)}^{r})}{\sum\nolimits_{i \in S_r} exp(h_{t-1}^{c} y_{i}^{r}) }</math><br />
: <math>P_{c}(w_t) = \frac{exp(h_{t}^{r} y_{c(w)}^{c})}{\sum\nolimits_{i \in S_c} exp(h_{t}^{r} y_{i}^{c}) }</math><br />
: <math> P(w_t) = P_{r}(w_t) P_{c}(w_t) </math> <br />
</center><br />
where <br />
<center><br />
:<math> r(w) </math> = row index of word w <br />
:<math> c(w) </math> = column index of word w<br />
:<math> y_{i}^{r} \in \mathbb{R}^m </math> = i-th vector of <math> Y^r \in \mathbb{R}^{m \times \sqrt{|V|}}</math> <br />
:<math> y_{i}^{c} \in \mathbb{R}^m </math> = i-th vector of <math> Y^c \in \mathbb{R}^{m \times \sqrt{|V|}}</math><br />
:<math> S_r </math> = the set of rows of the word table<br />
:<math> S_c </math> = the set of columns of the word table<br />
</center><br />
<br />
We can see that by using above equation, we effectively reduce the computation of the probability of the next word from a $|V|$-way normalization (in standard RNN models) to two $\sqrt {|V|}$-way normalizations. Note that we don't see the t-th word before predicting it. So in the above diagram, given the input column vector <math>x^c_{t-1} </math> of the (t-1)-th word, we first infer the row probability <math>P_r(w_t)</math> of the t-th word, and then choose the index of the row the largest probability in <math>P_r(w_t)</math> to look up the next input row vector <math>x^r_{t} </math>. Similarly, we can infer the column probability <math>P_c(w_t)</math> of the t-th word. <br />
<br />
Essentially, in LightRNN, the prediction of the word at time t (<math> w_t </math>) based on word at time t-1 (<math> w_{t-1} </math>) is achieved by selecting the index <math> r </math> and <math> c </math> with the highest probabilities <math> P_{r}(w_t) </math>, <math> P_{c}(w_t) </math>. Then, the probability of each word is computed based on the multiplication of <math> P_{r}(w_t) </math> and <math> P_{c}(w_t) </math>.<br />
<br />
== Part III: Bootstrap for Word Allocation ==<br />
<br />
As mentioned before, the major innovative aspect of LightRNN is the development of 2-component shared embedding. Such structure can be used in building a recurrent neural network called LightRNN. However, how should such word table representation be constructed is the key part of building a successful LightRNN model. In this section, the procedures of constructing 2C shared embedding structure is explained. <br />
The fundamental idea is using bootstrap method by minimizing a loss function (namely, negative log-likelihood function). The detailed procedures are described as the following:<br />
<br />
Step 1: First, all words in a vocabulary are randomly assigned to individual position within the word table<br />
<br />
Step 2: Train LightRNN model based on word table produced in step 1 until certain criteria are met<br />
<br />
Step 3: By fixing the training results of input and output embedding matrices (W & U) from step 2, adjust the position of words by minimizing the loss function over all the words. Then, repeat from step 2<br />
<br />
The authors presented the overall loss function for word w moving to position [i, j] using a negative log-likelihood function (NLL) as the following:<br />
<center><br />
<math> NLL = \sum\limits_{t=1}^T -logP(w_t) = \sum\limits_{t=1}^T -log[P_{r}(w_t) P_{c}(w_t)] = \sum\limits_{t=1}^T -log[P_{r}(w_t)] – log[P_{c}(w_t)] = \sum\limits_{w=1}^{|V|} NLL_w </math><br />
</center><br />
where <math> NLL_w </math> is the negative log-likelihood of a word w. <br />
<br />
Since in 2-component shared embedding structure, a word (w) is represented by one row vector and one column vector, <math> NLL_w </math> can be rewritten as <math> l(w, r(w), c(w)) </math> where <math> r(w) </math> and <math> c(w) </math> are the position index of word w in the word table. Next, the authors defined 2 more terms to explain the meaning of <math> NLL_w </math>: <math> l_r(w,r(w)) </math> and <math> l_c(w,c(w)) </math>, namely the row component and column component of <math> l(w, r(w), c(w)) </math>. The above can be summarised by the following formulas: <br />
<center><br />
<math> NLL_w = \sum\limits_{t \in S_w} -logP(w_t) = l(w, r(w), c(w)) </math> <br><br />
<math> = \sum\limits_{t \in S_w} -logP_r(w_t) + \sum\limits_{t \in S_w} -logP_c(w_t) = l_r(w,r(w)) + l_c(w,c(w))</math> <br><br />
<math> = \sum\limits_{t \in S_w} -log (\frac{exp(h_{t-1}^{c} y_{i}^{r})}{\sum\nolimits_{k} exp(h_{t-1}^{c} y_{i}^{k})}) + \sum\limits_{t \in S_w} -log (\frac{exp(h_{t}^{r} y_{j}^{c})}{\sum\nolimits_{k} exp(h_{t}^{r} y_{k}^{c}) }) </math> <br> <br />
where <math> S_w </math> is the set of all possible positions within the word table <br />
</center><br />
In summary, the overall loss function for word w to move to position [i, j] is the sum of its row loss and column loss of moving to position [i, j]. Therefore, total loss of moving to position [i, j] <math> l(w, i, j) = l_r(w, i) + l_c(w, j)</math>. Thus, to update the table by reallocating each word, we are looking for position [i, j] for each word w that minimize the total loss function, mathematically written as for the following:<br />
<center><br />
<math> \min\limits_{a} \sum\limits_{w,i,j} l(w,i,j)a(w,i,j) </math> such that <br><br />
<math> \sum\limits_{(i,j)} a(w,i,j) = 1 \space \forall w \in V, \sum\limits_{(w)} a(w,i,j) = 1 \space \forall i \in S_r, j \in S_j</math> <br><br />
<math> a(w,i,j) \in {0,1}, \forall w \in V, i \in S_r, j \in S_j</math> <br><br />
where <math> a(w,i,j) =1 </math> indicates moving word w to position [i, j]<br />
</center><br />
<br />
After calculating $l(w, i, j)$ for all possible $w, i, j$, the above optimization leads forcing $a(w, i, j)$ to be equal to 1 for $i, j$ in which $l(w, i, j)$ is minimum and 0 elsewhere (i.e. finding the best place for the word $w$ in the table).<br />
<br />
= LightRNN Example =<br />
<br />
After describing the theoretical background of the LightRNN algorithm, the authors applied this method to 2 datasets (2013 ACL Workshop Morphological Language Dataset (ACLW) & One-Billion-Word Benchmark Dataset (BillonW)) and compared its performance with several other state-of-the-art RNN algorithms. The following table shows some summary statistics of those 2 datasets:<br />
<br />
[[File:Table1YH.PNG|700px|thumb|centre|Table 1. Summary Statistics of Datasets]]<br />
<br />
The goal of a probabilistic language model is either to compute the probability distribution of a sequence of given words (e.g. <math> P(W) = P(w_1, w_2, … , w_n)</math>) or to compute the probability of the next word given some previous words (e.g. <math> P(w_5 | w_1, w_2, w_3, w_4)</math>) (Jurafsky, 2017). In this paper, the evaluation matrix for the performance of LightRNN algorithm is perplexity <math> PPL </math> which is defined as the following: <br />
<center><br />
<math> PPL = exp(\frac{NLL}{T})</math> <br><br />
where T = number of tokens in the test set<br />
</center><br />
<br />
Based on the mathematical definition of PPL, a well-performed model will have a lower perplexity. <br />
The authors then trained “LSTM-based LightRNN using stochastic gradient descent with truncated backpropagation through time” (Li, Qin, Yang, Hu, & Liu, 2016). To begin with, the authors first used the ACLW French dataset to determine the size of embedding matrix. From the results shown in Table 2, larger embedding size corresponds to higher accuracy rate (expressed in terms of perplexity). Therefore, they adopted embedding size of 1000 to be used in LightRNN to analyze the ACLW datasets. <br />
<br />
[[File:Table2YH.PNG|700px|thumb|centre|Table 2. Testing PPL of LightRNN on ACLW-French dataset w.r.t. embedding size]]<br />
<br />
'''Advantage 1: small model size'''<br />
<br />
One of the major advantages of using LightRNN on NLP tasks is significantly reduced model size, which means fewer number of parameters to estimate. By comparing LightRNN with two other RNN algorithms and the baseline language model with Kneser-Ney smoothing. Those two RNN algorithms are: HSM which uses LSTM RNN algorithm with hierarchical softmax for word prediction; C-HSM which uses both hierarchical softmax and character-level convolutional filters for input embedding. From the results table shown below, we can see that LightRNN has the lowest perplexity while keeping the model size significantly smaller compared to the other three algorithms. <br />
<br />
[[File:Table5YH.PNG|700px|thumb|centre|Table 3. PPL Results in test set on ACLW datasets]]<br />
Italic results are the previous state-of-the-art. #P denotes the number of parameters. <br />
<br />
'''Advantage 2: high training efficiency'''<br />
<br />
Another advantage of LightRNN model is its shorter training time while maintaining same level of perplexity compared to other RNN algorithms. When comparing to both C-HSM and HSM (shown below in Table 4), LightRNN only takes half the runtime but achieve same level of perplexity when applied to both ACLW and BillionW datasets. In the last column of Table 3, the amount of time used for word table reconstruction is presented as the percentage of the total runtime. As we can see, the training time for word reallocation takes up only a very small proportion of the total runtime. However, the resulting reconstructed word table can be used as a valuable output, which is further explained in the next section. <br />
<br />
[[File:Table3YH.PNG|700px|thumb|centre|Table 4. Runtime comparisons in order to achieve the HSMs’ baseline PPL]]<br />
<br />
<br />
'''Advantage 3: semantically valid word allocation table'''<br />
<br />
As explained in the previous section, LightRNN uses a word allocation table that gets updated in every iteration of the algorithm. The optimal structure of the table should assign semantically similar words onto the same row or column in order to reduce the number of parameters to estimate. Below is a snapshot of the reconstructed word table used in LightRNN algorithm. Evidently, we can see in row 887, all URL addresses are grouped together and in row 872 all verbs in past tense are grouped together. As the authors explained in the paper, LightRNN doesn’t assume independence of each word but instead using a shared embedding table. In this way, it reduces the model size by utilizing common embedding elements of the table/matrix, and also uses such preprocessed data to improve the efficiency of this algorithm.<br />
<br />
[[File:Table6YH.PNG|700px|thumb|centre|Table 6. Sample Word Allocation Table]]<br />
<br />
= Remarks =<br />
In summary, the proposed method in this paper is mainly on developing a new way of using word embedding. Words with similar semantic meanings are embedded using similar vectors. Those vectors are then divided into row and column components where similar words are grouped together by having shared row and column components in the word representation table. This is indeed a dimensional reduction, i.e. use the row and column "semantic vectors" to approximate the coded word. Because of this structural change of input word embedding, RNN model needs to adapt by having both row and column components being fed into the network. However, the fundamental structure of RNN model does not changed. Therefore, personally I would say it’s a new word embedding technique rather than a new development in model construction. One major confusion I have when reading this paper is how those row and column components in the word allocation table are determined. From the paper itself, the authors didn’t explain how they are constructed. <br />
<br />
Such shared word embedding technique are prevalently used in NLP. For instance, in language translation, similar words from different languages are grouped together so that the machine can translate sentences from one language to another. In Socher et al. (2013a), English and Chinese words are embedded in the same space so that we can find similar English (Chinese) words for Chinese (English) words. (Zou, Socher, Cer, & Manning, 2013). Word2vec is also a commonly used technique for word embedding, which uses a two-layer neural network to transform text into numeric vectors where similar words will have similar numeric values. The key feature of word2vec is that semantically similar words (which is now represented by numeric vectors) can be grouped together (“Word2vec,” n.d.; Bengio, Ducharme, & Vincent, 2001; Bengio, Ducharme, Vincent, & Jauvin, 2003).<br />
<br />
An interesting area of further exploration proposed by the authors is an extension of this method to k-component shared embeddings where k>2. Words probably share similar semantic meanings on more than two dimensions, and this extension could reduce network size even further. However it could also further complicate the bootstrapping phase of training.<br />
<br />
Code for LightRNN can be found on Github : <br />
Tensorflow : https://github.com/YisenWang/LightRNN-NIPS2016-Tensorflow_code<br />
= Reference =<br />
Bengio, Y, Ducharme, R., & Vincent, P. (2001). A Neural Probabilistic Language Model. In Journal of Machine Learning Research (Vol. 3, pp. 932–938). https://doi.org/10.1162/153244303322533223<br />
<br />
Bengio, Yoshua, Ducharme, R., Vincent, P., & Jauvin, C. (2003). A Neural Probabilistic Language Model. Journal of Machine Learning Research, 3(Feb), 1137–1155.<br />
<br />
Brownlee, J. (2017, September 20). 7 Applications of Deep Learning for Natural Language Processing. Retrieved October 27, 2017, from https://machinelearningmastery.com/applications-of-deep-learning-for-natural-language-processing/<br />
<br />
Hirschberg, J., & Manning, C. D. (2015). Advances in natural language processing. Science, 349(6245), 261–266. https://doi.org/10.1126/science.aaa8685<br />
<br />
Jurafsky, D. (2017, January). Language Modeling Introduction to N grams. Presented at the CS 124: From Languages to Information, Stanford University. Retrieved from https://web.stanford.edu/class/cs124/lec/languagemodeling.pdf<br />
<br />
Li, X., Qin, T., Yang, J., Hu, X., & Liu, T. (2016). LightRNN: Memory and Computation-Efficient Recurrent Neural Networks. Advances in Neural Information Processing Systems 29, 4385–4393.<br />
<br />
Recurrent Neural Networks. (n.d.). Retrieved October 8, 2017, from https://www.tensorflow.org/tutorials/recurrent<br />
<br />
Vector Representations of Words. (2017, August 17). Retrieved October 8, 2017, from https://www.tensorflow.org/tutorials/word2vec<br />
<br />
Word2vec. (n.d.). Retrieved October 26, 2017, from https://deeplearning4j.org/word2vec.html<br />
<br />
Zou, W. Y., Socher, R., Cer, D., & Manning, C. D. (2013). Bilingual word embeddings for phrase-based machine translation, 1393–1398.<br />
<br />
Kneser Ney Smoothing - : https://en.wikipedia.org/wiki/Kneser%E2%80%93Ney_smoothing & http://www.foldl.me/2014/kneser-ney-smoothing/<br />
<br />
Boné R., Assaad M., Crucianu M. (2003) Boosting Recurrent Neural Networks for Time Series Prediction. In: Pearson D.W., Steele N.C., Albrecht R.F. (eds) Artificial Neural Nets and Genetic Algorithms. Springer, Vienna</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=LightRNN:_Memory_and_Computation-Efficient_Recurrent_Neural_Networks&diff=30443LightRNN: Memory and Computation-Efficient Recurrent Neural Networks2017-11-15T18:21:28Z<p>Nkadiyar: /* LightRNN Structure */</p>
<hr />
<div>= Introduction =<br />
<br />
The study of natural language processing has been around for more than fifty years. It begins in 1950s which the specific field of natural language processing (NLP) is still embedded in the subject of linguistics (Hirschberg & Manning, 2015). After the emergence of strong computational power, computational linguistics began to evolve and gradually branch out to various applications in NLP, such as text classification, speech recognition and question answering (Brownlee, 2017). Computational linguistics or natural language processing is usually defined as “subfield of computer science concerned with using computational techniques to learn, understand, and produce human language content” (Hirschberg & Manning, 2015, p. 261). <br />
<br />
With the development of deep neural networks, one type of neural network, namely recurrent neural networks (RNN) have performed significantly well in many natural language processing tasks. The reason is that nature of RNN takes into account the past inputs as well as the current input without resulting in vanishing or exploding gradient. More detail of how RNN works in the context of NLP will be discussed in the section of NLP using RNN. However, one limitation of RNN used in NLP is its enormous size of input vocabulary. This will result in a very complex RNN model with too many parameters to train and makes the training process both time and memory-consuming. This serves as the major motivation for this paper’s authors to develop a new technique utilized in RNN, which is particularly efficient at processing large size of vocabulary in many NLP tasks, namely LightRNN.<br />
<br />
= Motivations =<br />
<br />
In language modelling, researchers used to represent words by arbitrary codes, such as “Id143” is the code for “dog” (“Vector Representations of Words,” 2017). Such coding of words is completely random, and it loses the meaning of the words and (more importantly) connection with other words. Nowadays, one-hot representation of words is commonly used, in which a word is represented by a vector of numbers and the dimension of the vector is related to the size of the vocabulary. In RNN, all words in the vocabulary are coded using one-hot representation and then mapped to an embedding vector (Li, Qin, Yang, Hu, & Liu, 2016). Such embedding vector is “a continuous vector space where semantically similar words are mapped to nearby points” (“Vector Representations of Words” 2017, para. 6). Popular RNN structure used in NLP task is long short-term memory (LSTM). In order to predict the probability of the next word, the last hidden layer of the network needs to calculate the probability distribution over all other words in the vocabulary. Note that the most time-consuming operation in RNNs is to calculate a probability distribution over all the words in the vocabulary, which requires the multiplication of the output-embedding matrix and the hidden state at each position of a sequence. Lastly, an activation function (commonly, softmax function) is used to select the next word with the highest probability. <br />
This method has 3 major limitations:<br />
<br />
1. Memory Constraint <br />
<br />
When input vocabulary contains an enormous amount of unique words, which is very common in various NLP tasks, the size of model becomes very large. This means the number of trainable parameters is very big, which makes it difficult to fit such model on a regular GPU device.<br />
<br />
2. Computationally Heavy to Train<br />
<br />
As previously mentioned, the probability distribution of all other words in the vocabulary needs to be computed to determine what predicted word it would be. When the size of the vocabulary is large, such calculation can be computationally heavy. <br />
<br />
3. Low Compressability<br />
<br />
Due to the memory and computation-consuming process of RNN applied in NLP tasks, mobile devices cannot usually handle such algorithm, which makes it undesirable and limits its usage.<br />
<br />
An alternative approach to handle the overhead is by leveraging weak lower-level learners by Boosting. But a drawback is that this technique has only been implemented for a few specific tasks in the past such as Time Series predictions [Boné et. al.].<br />
<br />
= LightRNN Structure =<br />
<br />
The authors of the paper proposed a new structure that effectively reduces the size of the model by arranging all words in the vocabulary into a word table, which is referred as “2-Component (2C) shared embedding for word representation”. This is done by factorizing a vocabulary's embedding into two shared components (row and column). Thus, a word is indexed by its location in such table, which in terms is characterized by the corresponding row and column components. Each row and column component are unique row vector and column vector respectively. By organizing each word in the vocabulary in this manner, multiple words can share the same row component or column component and it can reduce the number of trainable parameters significantly. <br />
The next question is how to construct such word table. More specifically, how to allocate each word in the vocabulary to different positions so that semantically similar words are in the same row or column. The authors proposed a bootstrap method to solve this problem. Essentially, we first randomly distribute words into the table. Then, we let the model “learn” better position of each word by minimizing training error. By repeating this process, each word can be allocated to a particular position within the table so that similar words share common row or column components. More details of those 2 parts of LightRNN structure will be discussed in the following sections.<br />
<br />
There are 2 major benefits of the proposed technique:<br />
<br />
1. Computationally efficient<br />
<br />
The name “LightRNN” is to illustrate the small model size and fast training speed. Because of these features of the new RNN architecture, it’s possible to launch such model onto regular GPU and other mobile devices. <br />
<br />
2. Higher scalability <br />
<br />
The authors briefly explained this algorithm is scalable because if parallel-computing is needed to train such model, the difficulty of combining smaller models is low. <br />
<br />
<br />
== Part I: 2-Component Shared Embedding ==<br />
<br />
The key aspect of LightRNN structure is its innovative method of word representation, namely 2-Component Shared Embedding. All words in the vocabulary are organized into a table with row components and column components. Each pair of the element in a row component and a column component is corresponding to a unique word in the vocabulary. For instance, the <math>i^{th}</math> row and <math>j^{th}</math> column are the row and column indexes for <math>X_{ij}</math>. As shown in the following graph, <math>x_{1}</math> is corresponding to the words “January”. In 2C shared embedding table, it’s indexed by 2 elements: <math>x^{r}_{1}</math> and <math>x^{c}_{1}</math> where the subscript indicates which row component and column component this word belongs to. Ideally, words that share similar semantic features should be assigned to the same row or column. The shared embedding word table in Figure 1 serves as a good example: the word “one” and “January” are assigned to the same column, while the word “one” and “two” are allocated to the same row. <br />
<br />
[[File:2C shared embedding.png|700px|thumb|centre|Fig 1. 2-Component Shared Embedding for Word Representation]]<br />
<br />
The main advantage of using such word representation is it reduces the number of vector/element needed for input word embedding. For instance, if there are 25 unique words in the vocabulary, the number of vectors to represent all the words is 10, namely 5 row vectors/elements and 5 column vectors/elements. Therefore, the shared embedding word table is a 5 by 5 matrix. In general, the formula for calculating number of vector/element needed to represent |V| words is <math>2\sqrt{|V|}</math>.<br />
<br />
== Part II: How 2C Shared Embedding is Used in LightRNN ==<br />
<br />
After constructing such word representation table, those 2-component shared embedding matrices are fed into the recurrent neural network. The following Figure 2 demonstrates a portion of LightRNN structure (left) with comparison with the regular RNN (right). Compared to regular RNN where a single input <math>x_{t-1}</math> is fed into the network each time, 2 elements of a single input <math>x_{t-1}</math>: <math>x^{r}_{t-1}</math> and <math>x^{c}_{t-1}</math> are fed into LightRNN. <br />
<br />
[[File:LightRNN.PNG |700px|thumb|centre|Fig 2. LightRNN Structure & Regular RNN]]<br />
<br />
As mentioned before, the last hidden layer will produce the probabilities of <math>word_{t}</math>. Based on the diagram below, the following formulas are used:<br />
Let n = dimension/length of a row input vector/a column input vector, <math>X^{c}, X^{r} \in \mathbb{R}^{n \times \sqrt{|V|}}</math> denotes the input-embedding matrices: <br />
<center><br />
: row vector <math>x^{r}_{t-1} \in \mathbb{R}^n</math><br />
: column vector <math>x^{c}_{t-1} \in \mathbb{R}^n</math><br />
</center><br />
<br />
Let <math>h^{c}_{t-1}, h^{r}_{t-1} \in \mathbb{R}^m</math> denotes the two hidden layers where m = dimension of the hidden layer:<br />
<center><br />
: <math>h^{c}_{t-1} = f(W x_{t-1}^{c} + U h_{t-1}^{r} + b) </math><br />
: <math>h^{r}_{t} = f(W x_{t}^{r} + U h_{t-1}^{c} + b) </math><br />
</center><br />
where <math>W \in \mathbb{R}^{m \times n}</math>, <math>U \in \mathbb{R}^{m \times m}</math>, and <math>b \in \mathbb{R}^m</math> and <math>f</math> is a nonlinear activation function<br />
<br />
The final step in LightRNN is to calculate <math>P_{r}(w_{t})</math> and <math>P_{c}(w_{t})</math> , which means the probability of a word w at time t, using the following formulas:<br />
<center><br />
: <math>P_{r}(w_t) = \frac{exp(h_{t-1}^{c} y_{r(w)}^{r})}{\sum\nolimits_{i \in S_r} exp(h_{t-1}^{c} y_{i}^{r}) }</math><br />
: <math>P_{c}(w_t) = \frac{exp(h_{t}^{r} y_{c(w)}^{c})}{\sum\nolimits_{i \in S_c} exp(h_{t}^{r} y_{i}^{c}) }</math><br />
: <math> P(w_t) = P_{r}(w_t) P_{c}(w_t) </math> <br />
</center><br />
where <br />
<center><br />
:<math> r(w) </math> = row index of word w <br />
:<math> c(w) </math> = column index of word w<br />
:<math> y_{i}^{r} \in \mathbb{R}^m </math> = i-th vector of <math> Y^r \in \mathbb{R}^{m \times \sqrt{|V|}}</math> <br />
:<math> y_{i}^{c} \in \mathbb{R}^m </math> = i-th vector of <math> Y^c \in \mathbb{R}^{m \times \sqrt{|V|}}</math><br />
:<math> S_r </math> = the set of rows of the word table<br />
:<math> S_c </math> = the set of columns of the word table<br />
</center><br />
<br />
We can see that by using above equation, we effectively reduce the computation of the probability of the next word from a $|V|$-way normalization (in standard RNN models) to two $\sqrt {|V|}$-way normalizations. Note that we don't see the t-th word before predicting it. So in the above diagram, given the input column vector <math>x^c_{t-1} </math> of the (t-1)-th word, we first infer the row probability <math>P_r(w_t)</math> of the t-th word, and then choose the index of the row the largest probability in <math>P_r(w_t)</math> to look up the next input row vector <math>x^r_{t} </math>. Similarly, we can infer the column probability <math>P_c(w_t)</math> of the t-th word. <br />
<br />
Essentially, in LightRNN, the prediction of the word at time t (<math> w_t </math>) based on word at time t-1 (<math> w_{t-1} </math>) is achieved by selecting the index <math> r </math> and <math> c </math> with the highest probabilities <math> P_{r}(w_t) </math>, <math> P_{c}(w_t) </math>. Then, the probability of each word is computed based on the multiplication of <math> P_{r}(w_t) </math> and <math> P_{c}(w_t) </math>.<br />
<br />
== Part III: Bootstrap for Word Allocation ==<br />
<br />
As mentioned before, the major innovative aspect of LightRNN is the development of 2-component shared embedding. Such structure can be used in building a recurrent neural network called LightRNN. However, how should such word table representation be constructed is the key part of building a successful LightRNN model. In this section, the procedures of constructing 2C shared embedding structure is explained. <br />
The fundamental idea is using bootstrap method by minimizing a loss function (namely, negative log-likelihood function). The detailed procedures are described as the following:<br />
<br />
Step 1: First, all words in a vocabulary are randomly assigned to individual position within the word table<br />
<br />
Step 2: Train LightRNN model based on word table produced in step 1 until certain criteria are met<br />
<br />
Step 3: By fixing the training results of input and output embedding matrices (W & U) from step 2, adjust the position of words by minimizing the loss function over all the words. Then, repeat from step 2<br />
<br />
The authors presented the overall loss function for word w moving to position [i, j] using a negative log-likelihood function (NLL) as the following:<br />
<center><br />
<math> NLL = \sum\limits_{t=1}^T -logP(w_t) = \sum\limits_{t=1}^T -log[P_{r}(w_t) P_{c}(w_t)] = \sum\limits_{t=1}^T -log[P_{r}(w_t)] – log[P_{c}(w_t)] = \sum\limits_{w=1}^{|V|} NLL_w </math><br />
</center><br />
where <math> NLL_w </math> is the negative log-likelihood of a word w. <br />
<br />
Since in 2-component shared embedding structure, a word (w) is represented by one row vector and one column vector, <math> NLL_w </math> can be rewritten as <math> l(w, r(w), c(w)) </math> where <math> r(w) </math> and <math> c(w) </math> are the position index of word w in the word table. Next, the authors defined 2 more terms to explain the meaning of <math> NLL_w </math>: <math> l_r(w,r(w)) </math> and <math> l_c(w,c(w)) </math>, namely the row component and column component of <math> l(w, r(w), c(w)) </math>. The above can be summarised by the following formulas: <br />
<center><br />
<math> NLL_w = \sum\limits_{t \in S_w} -logP(w_t) = l(w, r(w), c(w)) </math> <br><br />
<math> = \sum\limits_{t \in S_w} -logP_r(w_t) + \sum\limits_{t \in S_w} -logP_c(w_t) = l_r(w,r(w)) + l_c(w,c(w))</math> <br><br />
<math> = \sum\limits_{t \in S_w} -log (\frac{exp(h_{t-1}^{c} y_{i}^{r})}{\sum\nolimits_{k} exp(h_{t-1}^{c} y_{i}^{k})}) + \sum\limits_{t \in S_w} -log (\frac{exp(h_{t}^{r} y_{j}^{c})}{\sum\nolimits_{k} exp(h_{t}^{r} y_{k}^{c}) }) </math> <br> <br />
where <math> S_w </math> is the set of all possible positions within the word table <br />
</center><br />
In summary, the overall loss function for word w to move to position [i, j] is the sum of its row loss and column loss of moving to position [i, j]. Therefore, total loss of moving to position [i, j] <math> l(w, i, j) = l_r(w, i) + l_c(w, j)</math>. Thus, to update the table by reallocating each word, we are looking for position [i, j] for each word w that minimize the total loss function, mathematically written as for the following:<br />
<center><br />
<math> \min\limits_{a} \sum\limits_{w,i,j} l(w,i,j)a(w,i,j) </math> such that <br><br />
<math> \sum\limits_{(i,j)} a(w,i,j) = 1 \space \forall w \in V, \sum\limits_{(w)} a(w,i,j) = 1 \space \forall i \in S_r, j \in S_j</math> <br><br />
<math> a(w,i,j) \in {0,1}, \forall w \in V, i \in S_r, j \in S_j</math> <br><br />
where <math> a(w,i,j) =1 </math> indicates moving word w to position [i, j]<br />
</center><br />
<br />
After calculating $l(w, i, j)$ for all possible $w, i, j$, the above optimization leads forcing $a(w, i, j)$ to be equal to 1 for $i, j$ in which $l(w, i, j)$ is minimum and 0 elsewhere (i.e. finding the best place for the word $w$ in the table).<br />
<br />
= LightRNN Example =<br />
<br />
After describing the theoretical background of the LightRNN algorithm, the authors applied this method to 2 datasets (2013 ACL Workshop Morphological Language Dataset (ACLW) & One-Billion-Word Benchmark Dataset (BillonW)) and compared its performance with several other state-of-the-art RNN algorithms. The following table shows some summary statistics of those 2 datasets:<br />
<br />
[[File:Table1YH.PNG|700px|thumb|centre|Table 1. Summary Statistics of Datasets]]<br />
<br />
The goal of a probabilistic language model is either to compute the probability distribution of a sequence of given words (e.g. <math> P(W) = P(w_1, w_2, … , w_n)</math>) or to compute the probability of the next word given some previous words (e.g. <math> P(w_5 | w_1, w_2, w_3, w_4)</math>) (Jurafsky, 2017). In this paper, the evaluation matrix for the performance of LightRNN algorithm is perplexity <math> PPL </math> which is defined as the following: <br />
<center><br />
<math> PPL = exp(\frac{NLL}{T})</math> <br><br />
where T = number of tokens in the test set<br />
</center><br />
<br />
Based on the mathematical definition of PPL, a well-performed model will have a lower perplexity. <br />
The authors then trained “LSTM-based LightRNN using stochastic gradient descent with truncated backpropagation through time” (Li, Qin, Yang, Hu, & Liu, 2016). To begin with, the authors first used the ACLW French dataset to determine the size of embedding matrix. From the results shown in Table 2, larger embedding size corresponds to higher accuracy rate (expressed in terms of perplexity). Therefore, they adopted embedding size of 1000 to be used in LightRNN to analyze the ACLW datasets. <br />
<br />
[[File:Table2YH.PNG|700px|thumb|centre|Table 2. Testing PPL of LightRNN on ACLW-French dataset w.r.t. embedding size]]<br />
<br />
'''Advantage 1: small model size'''<br />
<br />
One of the major advantages of using LightRNN on NLP tasks is significantly reduced model size, which means fewer number of parameters to estimate. By comparing LightRNN with two other RNN algorithms and the baseline language model with Kneser-Ney smoothing. Those two RNN algorithms are: HSM which uses LSTM RNN algorithm with hierarchical softmax for word prediction; C-HSM which uses both hierarchical softmax and character-level convolutional filters for input embedding. From the results table shown below, we can see that LightRNN has the lowest perplexity while keeping the model size significantly smaller compared to the other three algorithms. <br />
<br />
[[File:Table5YH.PNG|700px|thumb|centre|Table 3. PPL Results in test set on ACLW datasets]]<br />
Italic results are the previous state-of-the-art. #P denotes the number of parameters. <br />
<br />
'''Advantage 2: high training efficiency'''<br />
<br />
Another advantage of LightRNN model is its shorter training time while maintaining same level of perplexity compared to other RNN algorithms. When comparing to both C-HSM and HSM (shown below in Table 4), LightRNN only takes half the runtime but achieve same level of perplexity when applied to both ACLW and BillionW datasets. In the last column of Table 3, the amount of time used for word table reconstruction is presented as the percentage of the total runtime. As we can see, the training time for word reallocation takes up only a very small proportion of the total runtime. However, the resulting reconstructed word table can be used as a valuable output, which is further explained in the next section. <br />
<br />
[[File:Table3YH.PNG|700px|thumb|centre|Table 4. Runtime comparisons in order to achieve the HSMs’ baseline PPL]]<br />
<br />
<br />
'''Advantage 3: semantically valid word allocation table'''<br />
<br />
As explained in the previous section, LightRNN uses a word allocation table that gets updated in every iteration of the algorithm. The optimal structure of the table should assign semantically similar words onto the same row or column in order to reduce the number of parameters to estimate. Below is a snapshot of the reconstructed word table used in LightRNN algorithm. Evidently, we can see in row 887, all URL addresses are grouped together and in row 872 all verbs in past tense are grouped together. As the authors explained in the paper, LightRNN doesn’t assume independence of each word but instead using a shared embedding table. In this way, it reduces the model size by utilizing common embedding elements of the table/matrix, and also uses such preprocessed data to improve the efficiency of this algorithm.<br />
<br />
[[File:Table6YH.PNG|700px|thumb|centre|Table 6. Sample Word Allocation Table]]<br />
<br />
= Remarks =<br />
In summary, the proposed method in this paper is mainly on developing a new way of using word embedding. Words with similar semantic meanings are embedded using similar vectors. Those vectors are then divided into row and column components where similar words are grouped together by having shared row and column components in the word representation table. This is indeed a dimensional reduction, i.e. use the row and column "semantic vectors" to approximate the coded word. Because of this structural change of input word embedding, RNN model needs to adapt by having both row and column components being fed into the network. However, the fundamental structure of RNN model does not changed. Therefore, personally I would say it’s a new word embedding technique rather than a new development in model construction. One major confusion I have when reading this paper is how those row and column components in the word allocation table are determined. From the paper itself, the authors didn’t explain how they are constructed. <br />
<br />
Such shared word embedding technique are prevalently used in NLP. For instance, in language translation, similar words from different languages are grouped together so that the machine can translate sentences from one language to another. In Socher et al. (2013a), English and Chinese words are embedded in the same space so that we can find similar English (Chinese) words for Chinese (English) words. (Zou, Socher, Cer, & Manning, 2013). Word2vec is also a commonly used technique for word embedding, which uses a two-layer neural network to transform text into numeric vectors where similar words will have similar numeric values. The key feature of word2vec is that semantically similar words (which is now represented by numeric vectors) can be grouped together (“Word2vec,” n.d.; Bengio, Ducharme, & Vincent, 2001; Bengio, Ducharme, Vincent, & Jauvin, 2003).<br />
<br />
An interesting area of further exploration proposed by the authors is an extension of this method to k-component shared embeddings where k>2. Words probably share similar semantic meanings on more than two dimensions, and this extension could reduce network size even further. However it could also further complicate the bootstrapping phase of training.<br />
<br />
Code for LightRNN can be found on Github : <br />
Tensorflow : https://github.com/YisenWang/LightRNN-NIPS2016-Tensorflow_code<br />
= Reference =<br />
Bengio, Y, Ducharme, R., & Vincent, P. (2001). A Neural Probabilistic Language Model. In Journal of Machine Learning Research (Vol. 3, pp. 932–938). https://doi.org/10.1162/153244303322533223<br />
<br />
Bengio, Yoshua, Ducharme, R., Vincent, P., & Jauvin, C. (2003). A Neural Probabilistic Language Model. Journal of Machine Learning Research, 3(Feb), 1137–1155.<br />
<br />
Brownlee, J. (2017, September 20). 7 Applications of Deep Learning for Natural Language Processing. Retrieved October 27, 2017, from https://machinelearningmastery.com/applications-of-deep-learning-for-natural-language-processing/<br />
<br />
Hirschberg, J., & Manning, C. D. (2015). Advances in natural language processing. Science, 349(6245), 261–266. https://doi.org/10.1126/science.aaa8685<br />
<br />
Jurafsky, D. (2017, January). Language Modeling Introduction to N grams. Presented at the CS 124: From Languages to Information, Stanford University. Retrieved from https://web.stanford.edu/class/cs124/lec/languagemodeling.pdf<br />
<br />
Li, X., Qin, T., Yang, J., Hu, X., & Liu, T. (2016). LightRNN: Memory and Computation-Efficient Recurrent Neural Networks. Advances in Neural Information Processing Systems 29, 4385–4393.<br />
<br />
Recurrent Neural Networks. (n.d.). Retrieved October 8, 2017, from https://www.tensorflow.org/tutorials/recurrent<br />
<br />
Vector Representations of Words. (2017, August 17). Retrieved October 8, 2017, from https://www.tensorflow.org/tutorials/word2vec<br />
<br />
Word2vec. (n.d.). Retrieved October 26, 2017, from https://deeplearning4j.org/word2vec.html<br />
<br />
Zou, W. Y., Socher, R., Cer, D., & Manning, C. D. (2013). Bilingual word embeddings for phrase-based machine translation, 1393–1398.<br />
<br />
Kneser Ney Smoothing - : https://en.wikipedia.org/wiki/Kneser%E2%80%93Ney_smoothing & http://www.foldl.me/2014/kneser-ney-smoothing/<br />
<br />
Boné R., Assaad M., Crucianu M. (2003) Boosting Recurrent Neural Networks for Time Series Prediction. In: Pearson D.W., Steele N.C., Albrecht R.F. (eds) Artificial Neural Nets and Genetic Algorithms. Springer, Vienna</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=%22Why_Should_I_Trust_You%3F%22:_Explaining_the_Predictions_of_Any_Classifier&diff=30442"Why Should I Trust You?": Explaining the Predictions of Any Classifier2017-11-15T18:15:57Z<p>Nkadiyar: /* Fidelity-Interpretability Trade-Off */</p>
<hr />
<div>==Introduction==<br />
<br />
Understanding why machine learning models behave the way they do helps users in model selection, feature engineering, and finally trusting the model to deploy it. In many cases even though a non-interpretable model is more accurate on the validation data set than the interpretable one, an interpretable model is chosen. But restricting the models to just an interpretable model isn't the best option. In this paper, the authors argue for explaining machine learning predictions using model-agnostic approaches and propose LIME (Local Interpretable Model-Agnostic Explanations), a novel explanation technique that explains the predictions of any classifier in an interpretable and faithful manner, by learning an interpretable model locally around the prediction. They also propose a method called sub-modular optimization to explain models globally by selecting a representative individual predictions. <br />
<br />
Some recent work aims to anticipate failures in machine learning, specifically for vision tasks. Failure modes are found by grouping incorrectly classified images according to a semantic attribute space, which are descriptions explaining the cause of the failure [12]. Another approach is to train a model that can predict failure [14]. Letting users know when the systems are likely to fail can lead to an increase in trust, by avoiding “silly mistakes” [13]. These solutions either require additional annotations and feature engineering that is specific to vision tasks or do not provide insight into why a decision should not be trusted. Furthermore, they assume that the current evaluation metrics are reliable, which may not be the case if problems such as data leakage are present.<br />
<br />
The experiments are conducted mainly on text data and image data, where models used are random forest (RF) and Convolutional Neural Networks (CNN). The experiments highlight utility of explanations in deciding between the models, which model is to be trusted and which is not and also improving the model based on the explanations. The main contributions of this paper are summarized as follows.<br />
<br />
* LIME, an algorithm to explain any prediction of any model (classifier or regressor) by approximating locally with an interpretable model that is faithful to the original model locally. <br />
<br />
* SP-LIME, a method to select final set of representative samples using sub-modular optimization to show to the user that pretty much captures what the original model is doing globally. <br />
<br />
* A detailed set of experiments with simulated subjects to prove the usefulness of LIME and SP-LIME. <br />
<br />
[[File:LIME.jpg|thumb|550px|Figure 1: Explaining individual predictions. A model predicts that a patient has the flu, and LIME highlights the symptoms in the patient’s history that led to the prediction. Sneeze and headache are portrayed as contributing to the “flu” prediction, while “no fatigue” is evidence against it. With these, a doctor can make an informed decision about whether to trust the model’s prediction.]]<br />
<br />
== Need for Explanations ==<br />
<br />
Prediction explanation basically means answering question such as "What triggers the model to make such a prediction?", and an answer to that would be to say, these are the features with certain range of values in your input sample that are contributing to the prediction the most. For example, it can be words in a document for text classification or patches of pixels in an image for image classification. And these explanations need to make sense to the user and hence has to be simple. Figure 1 illustrates an explanation procedure. In this case, an explanation is a small weighted list symptoms that either contribute to the prediction (in green) or are evidence against it (in red). It is very clear that the life of a doctor is much easier in terms of making a decision with the help of a model if intelligible explanations are provided. This is similar to the topic of inference versus prediction. Some people tend to think of these two terms are the same. However, statistically, their meanings are very different. Inference is given some datasets, we would like to infer how the output/response is generated based on a sequence of explanatory variables. For instance, we would like to know/infer how education level affect people’s income level. Prediction , on the other hand is by using a given dataset, we fit/train a model that will correctly predict the outcome of a new observation. <br />
<br />
Even with highest accuracy on validation dataset, we sometimes can't judge how the model is going to behave on other dataset. There are many reasons why this can happen. For example, Data leakage where some of the features used for training are heavily correlated with the target value in both training and validation data that might result in great train and validation results but will be of no use if used on all together a new dataset. But if explanations such as the one in Figure 1 are given, it becomes easy to fix the fault by removing the heavily correlated feature from the data and train. This is how you convert an untrustworthy model to a trustworthy one. Another problem is dataset shift, where train and test data come from different distributions and providing explanations for predictions will help user to take measures to make model generalize better.<br />
<br />
Figure 2 shows how explanations help in selecting between the models in addition to accuracy measures. In this case, by looking at the explanations it is easy to say that model with higher accuracy is in-fact the worst. Further, many a times, there can be difference between what we want the model to optimize and what the model is actually optimizing which might result into model not behaving as well as expected (For example in binary classification using logistic regression, all we are concerned is with accuracy, but the model is trying to reduce the log-loss due to the learning algorithms limitations). While we may not be able to quantitatively measure what difference has that made, but we still have some intuitions as to how the model has to behave with respect to some of its features, and if explanations are given for model predictions, we will be able to ascertain if those explanations make any sense. Also explaining models being model agnostic helps as we then compare different class of models in making a choice.<br />
<br />
[[File:explanation_example.png|thumb|550px|Figure 2: Explaining individual predictions of competing classifiers trying to determine if a document is about “Christianity” or “Atheism”. The bar chart represents the importance given to the most relevant words, also highlighted in the text. Color indicates which class the word contributes to (green for “Christianity”, magenta for “Atheism”).]]<br />
<br />
===Desired Characteristics for Explainers===<br />
<br />
Explainers have to be simple models so that humans are able to comprehend what the model is doing. For example, even with linear models, if the number of features are too many, it becomes hard for humans to comprehend overwhelming set of features. In case of models trained with word-embedding as features, we can't give these word-embeddings as explanations, instead they need to something different than these features.<br />
<br />
Another requirement of a good explainer is that it has to be locally faithful to the original model, i.e, it must behave similar to the model in the vicinity of instance being predicted. And these local explanations should aid in forming global explanations.<br />
<br />
==Local Interpretable Model-Agnostic Explanations (LIME)==<br />
<br />
The overall goal of LIME is to basically identify an interpretable model over the interpretable representation (features) that is locally faithful to the classifier. For example in case of text classification, a possible interpretable representation of data would be to use bag of words (one-hot encodings) instead of the word embeddings. Similarly for image classification, it may be a binary vector (one-hot encoding) indicating the “presence” or “absence” of a contiguous patch of pixels (a super-pixel or a segment), while the classifier may represent the image as a tensor with three color channels per pixel. Let $x ∈ R^d$ be the original representation of an instance being explained, and $x' ∈ \{0, 1\}^{d'}$ be a binary vector for its interpretable representation.<br />
<br />
=== Fidelity-Interpretability Trade-Off ===<br />
Formally, the authors define an explanation as a model $g ∈ G$, where $G$ is a class of potentially interpretable models, such as linear models, decision trees, or falling rule lists [3]. The domain of $g$ is $\{0, 1\}^{d'}$, i.e. $g$ acts over absence/presence of the interpretable components. Let $Ω(g)$ be a measure of complexity (like a regularizer term) of the explanation $g ∈ G$. For example, for decision trees $Ω(g)$ may be the depth of the tree, for linear models, $Ω(g)$ may be the number of non-zero weights. Let the model being explained be denoted $f : R^d → R$. In classification, $f(x)$ is the probability that $x$ belongs to a certain class. We further use $π_x(z)$ as a proximity measure between an instance $z$ to $x$, so as to define locality around $x$. Finally, let $L(f, g, π_x)$ be a measure of error between $g$ and $f$ in the locality defined by $π_x$. Our task is to minimize $L(f, g, π_x)$ while having $Ω(g)$ be low enough to be interpretable by humans. The minimum value of $L(f, g, π_x)$ provides an approximate balance measure for preserving interpretability and local fidelity(it should correspond to how the model behaves in the vicinity of the instance being predicted). The explanation produced by LIME is obtained by the following:<br />
<br />
:<math>\xi(x) = \underset{g\in\mathbb{G}}{\operatorname{argmin}}\, \mathcal{L}(f,g,\pi_x) + \Omega(g)\,\,\,\,\,\,\,\,\,\,\,\,\,\,\, Eq. (1)</math><br />
<br />
Even though $G$ can be any class of interpretable models, in this paper only sparse linear models are considered.<br />
<br />
=== Sampling for Local Exploration ===<br />
<br />
In order to learn the local behavior of $f$ using $g$ as the interpretable inputs vary, $L(f, g, π_x)$ is approximated by drawing samples, weighted by $π_x$. We sample instances around $x'$ by setting few features of $x'$ to 0, uniformly at random (where the number of such features set to 0 is also uniformly sampled). Given a perturbed sample $z' ∈ \{0, 1\}^{d'}$ (which contains a fraction of the features of $x'$ set to 0), we recover the sample in the original representation $z ∈ R^d$ and obtain the prediction $f(z)$ from the original classifier, and is used as a label for the explanation model. Given this dataset $Z$ of perturbed samples with the associated labels, we optimize Eq. (1) to get an explanation $ξ(x)$. The primary intuition behind LIME is presented in Figure 3, where we sample instances both in the vicinity of $x$ (which have a high weight due to $π_x$) and far away from $x$ (low weight from $π_x$). Even though the original model may be too complex to explain globally, LIME presents an explanation that is locally faithful (linear in this case), where the locality is captured by weight factor, $π_x$. A concrete example of this process is presented in the next section.<br />
<br />
[[File:decision_boundary.png|thumb|Figure 3: Toy example to present intuition for LIME. The black-box model’s complex decision function $f$ (unknown to LIME) is represented by the blue/pink background, which cannot be approximated well by a linear model. The bold red cross is the instance being explained. LIME samples instances, gets predictions using $f$, and weighs them by the proximity to the instance being explained (represented here by size). The dashed line is the learned explanation that is locally (but not globally) faithful.]]<br />
<br />
===Sparse Linear Explanations===<br />
<br />
For the rest of this paper, we let $G$ be the class of linear models, such that $g(z') = w_g · z'$ . We use the locally weighted square loss as $L$, as defined in Eq. (2), where we let $π_x(z) = exp(−D(x, z)^2 /σ^2 )$ be an exponential kernel defined on some distance function $D$ (e.g. cosine distance for text, L2 distance for images) with width $σ$ resulting in higher sample weights for perturbed samples that are in the vicinity of sample being explained and lower weights for samples that are far away.<br />
<br />
:<math>\mathcal{L}(f,g,\pi_x) = \underset{z,z'\in\mathbb{Z}}{\operatorname{\Sigma}} \pi_x(z)\,(f(z) - g(z'))\,\,\,\,\,\,\,\,\,\,\,\,\,\,\, Eq. (2)</math><br />
<br />
For text classification, the features used for explanation models are bag of words, and we set a limit on the number of words used for explanation, $K$. i.e. $Ω(g) = ∞1[||w_g||_0 > K]$. We use the same $Ω$ for image classification, using “super-pixels” (computed using any standard image segmentation algorithm, eg. quickshift() function in skimage python library [5]) instead of words, such that the interpretable representation of an image is a binary vector where 1 indicates the original super-pixel and 0 indicates a grayed out super-pixel. This particular choice of $Ω$ makes directly solving Eq. (1) intractable, but we approximate it by first selecting $K$ features with Lasso (using the regularization path [4]) (scikit lars_path method [6]) and then learning the weights via least squares (together a procedure we call K-LASSO in Algorithm 1). . <br />
<br />
There are some limitations to class of locally interpretable models $G$. It can happen that no $g \in G$ is powerful enough to be locally faithful to the original model, i.e, $g$ has a high bias w.r.t $f$ locally. This may be because the underlying model is quite complex even locally. Another issue could be, the features used for building explanations are not representative of the inputs or factors that the underlying model relies on. For example, a model that predicts sepia-toned images to be retro cannot be explained by presence of absence of super pixels.<br />
<br />
<br />
[[File:algorithm1.png|thumb|550px]]<br />
<br />
==Examples==<br />
<br />
===Example 1: Text classification with SVMs===<br />
<br />
In Figure 2 (RHS), Here LIME explains an SVM classifier with RBF kernel trained on "20 newsgroup dataset" to classify documents into "Atheism" and "Christianity". Even though it achieves 94% accuracy on validation dataset, the explanation for an instance show that features that are used the most in predictions are very arbitrary such as words like "Posting", "Host" and "Re". Even if these stop words or headers are removed, the classier still considers the proper names of people writing the post to be important features, which clearly doesn't make sense. Hence we can conclude by looking at the explanations that the dataset has serious issues and when a classifier is trained on it, it wont be able to generalize well. Hence suitable steps needed to be taken to train a trustworthy classifier.<br />
<br />
===Example 2: Deep networks for images===<br />
<br />
In this example we use sparse linear explanations to explain Google’s pre-trained Inception neural network [7] which is an image classifier. Here features used for giving explanations are super-pixels (segments) which intuitively makes sense as we humans detect an object in an image based on certain parts of the object. Figure 4a show an arbitrary image that we want to explain. Figures 4b, 4c, 4d show the super-pixels that contribute the most while predicting the image to be one of the top 3 classes. Here we set max feature count for explainer, $K = 10$. It is quite intuitive from Figure 4b in particular that why acoustic guitar was predicted to be electric based on the fretboard that resembles acoustic guitar. So even though the prediction is incorrect, it still is not unreasonable at all.<br />
<br />
[[File:inception_example.png|thumb|center|550px]]<br />
<br />
==Submodular Pick for Explaining Models==<br />
<br />
Submodular pick provides some sort of a global understanding of the model by methodically picking a set on non-redundant sample explanations. This method is still model agnostic and complimentary to calculating validation accuracy in machine learning problem. If we choose a large number of instances to explain, the user may not be able to go through each of them and verify what is wrong or right with the model. We say humans have a budget $B$, that is the number of instance explanations that they are willing to examine. So given a set of $X$ instances, the task of sub-modular optimization is to pick a maximum of $B$ instances that effectively capture the essence of what the model is doing in a non-redundant way.<br />
<br />
Here we have a explanations for $n$ instances and we construct an $n × d'$ matrix $W$, each row of which represent local importances of features used for each instance explanations. In case of using sparse linear models as explanations, for an instance $x_i$ and explanation $g_i = ξ(x_i)$, we set $W_{ij} = |w_{gij}|$. Further, for each feature (column) $j$ in $W$, let $I_j$ denote the global importance of that component in the explanation space. Intuitively, $I$ represents global importance of each feature. In Figure 5, we show a toy example $W$, with $n = d' = 5$, where $W$ is binary (for simplicity). Since feature $f2$ is used in most of the explanations, $I$ should have higher value for feature $f2$ than $f1$, i.e. $I_2 > I_1$. Formally for the text classifiers, we set $I_j = \sqrt{\sum_{i=1}^nW_{ij}}$ . For images, $I$ must represent something that can be compared across the super-pixels (segments) in different images, such as color histograms or other features of super-pixels; But this is left as a future work for now.<br />
<br />
While picking the samples, we must be careful in not picking up samples that explain the same thing. So we want minimum set of samples that cover maximum number of features as possible. In Figure 5, after the second row is picked, the third row is redundant, as the user has already seen features $f2$ and $f3$ - while the last row gives the user completely new features. Hence selecting the second and last row results in giving the maximum coverage of features. Eq. (3) formalizes this process, where we define coverage as the set function $c$ that, given $W$ and $I$, computes the total importance of the features that at least appear in one instance in a set $V$ .<br />
<br />
:<math>c(V, W, I) = {{\sum_{j = 1}^{d'}}}\mathbb{1}_{[\exists i\in V:W_{ij}\gt 0]}I_j\,\,\,\,\,\,\,\,\,\,\,\,\,\,\, Eq. (3)</math><br />
<br />
The pick problem, defined in Eq. (4), consists of finding the set $V, |V| ≤ B$ that achieves highest coverage.<br />
<br />
:<math>Pick(W, I) = \underset{V,|V|\leq B}{\operatorname{argmax}}\, c(V,W, I)\,\,\,\,\,\,\,\,\,\,\,\,\,\,\, Eq. (4)</math><br />
<br />
The problem in Eq. (4) is maximizing a weighted coverage function, and is NP-hard [10]. Let $c(V\cup\{i\}, W, I)−c(V, W, I)$ be the marginal coverage gain of adding an instance $i$ to a set $V$ . Due to sub-modularity, a greedy algorithm that iteratively adds the instance with the highest marginal coverage gain to the solution offers a constant-factor approximation guarantee of $1−1/e$ to the optimum [8]. This approximation process is outlined in Algorithm 2 and is called sub-modular pick.<br />
<br />
[[File:algorithm2.png|thumb|550px]] <br />
[[File:toy_example.png|thumb|Figure 5: Toy example W. Rows represent instances (documents) and columns represent features (words). Feature f2 (dotted blue) has the highest importance. Rows 2 and 5 (in red) would be selected by the pick procedure, covering all but feature f1.]]<br />
<br />
==Simulated User Experiments==<br />
<br />
This section highlights simulated user experiments conducted to evaluate how explanations help the user in different tasks. In particular, following questions are addressed: (1) Are the explanations faithful to the model, that is are the explainers correctly highlighting the important features that the model is selecting at least locally? (2) Can the explanations help users in ascertaining trust in predictions that is to say if the predictions make sense or not and (3) Are the explanations useful for evaluating the model as a whole and thereby aid in finally selecting a model. Code and data for experiments are available at https://github.com/marcotcr/lime-experiments.<br />
<br />
===Experiment Setup===<br />
<br />
2 sentiment analysis datasets (books and DVDs, 2000 instances each) where the task is to classify product reviews as positive or negative [9] are used in the experiment. Decision trees (DT), logistic regression with L2 regularization (LR), nearest neighbors (NN), and SVMs with RBF kernel, all using bag of words as features are trained separately. There is also a random forest (with 1000 trees) trained over average word2vec embeddings (RF). Unless otherwise noted, default methods of scikit are used. Dataset is divided into train and test in the ratio of 4:1. To explain individual predictions, LIME is compared with parzen [10], a method that approximates the black box classifier globally with Parzen windows, and explains individual predictions by taking the gradient of the prediction probability function. (Parzen-window density estimation is essentially a data-interpolation technique. Given an instance of the random sample, ${\bf x}$, Parzen-windowing estimates the PDF $P(X)$ from which the sample was derived. It essentially superposes kernel functions placed at each observation or datum. In this way, each observation $x_i$ contributes to the PDF estimate [15]). For parzen, $K$ features are selected as explanations such that they have the highest absolute gradients. Hyper-parameters for LIME and parzen are set using cross-validations and number of perturbed samples for a single explanation, $N$ is set to 15000. There is also a comparison made against a greedy procedure in which features that contribute the most to the predicted class are greedily removed until prediction changes or we reach a maximum of $K$ features that are removed, and these $K$ or less features are finally selected to explain the model. Also also a random procedure that randomly picks K features as an explanation is also used for comparison. For all the experiments below $K$ is set to 10. For experiments where the pick procedure applies, it's either a random pick (RP) or sub-modular pick (SP). We refer to pick-explainer combinations by adding RP or SP as a prefix.<br />
<br />
[[File:recall1.png|thumb|Figure 6: Recall [11] on truly important features for two interpretable classifiers on the books dataset.]] [[File:recall2.png|thumb|Figure 7: Recall on truly important features for two interpretable classifiers on the DVDs dataset.]]<br />
<br />
===Are explanations faithful to the model?===<br />
<br />
Here we evaluate faithfulness of explanations on classifiers that are by themselves explicable such as sparse logistic regression or decision trees. So we train 2 models, a sparse logistic regressor and a decision tree such that maximum number of features they use 10 so we now know what the important features (gold features) really are for these 2 models. For each prediction on the test set, we generate explanations and compute the fraction of these gold features that are recovered by the explanations (recall). We report this recall averaged over all the test instances in Figures 6 and 7. From the bar graphs we can see that the greedy approach is comparable to parzen on logistic regression, but is substantially worse on decision trees since changing a single feature at a time often does not have an effect on the prediction. The overall recall by parzen is low, likely due to the difficulty in approximating the original high-dimensional classifier. LIME consistently provides > 90% recall for both classifiers on both datasets, and this demonstrates that LIME explanations are more faithful to the model than any other explanations.<br />
<br />
===Should I trust this prediction?===<br />
<br />
"""In order to simulate trust in individual predictions, we first randomly select 25% of the features to be “untrustworthy”, and assume that the users can identify and would not want to trust these features (such as the headers in 20 newsgroups, leaked data, etc). We thus develop oracle “trustworthiness” by labeling test set predictions from a black box classifier as “untrustworthy” if the prediction changes when untrustworthy features are removed from the instance, and “trustworthy” otherwise. In order to simulate users, we assume that users deem predictions untrustworthy from LIME and parzen explanations if the prediction from the linear approximation changes when all untrustworthy features that appear in the explanations are removed (the simulated human “discounts” the effect of untrustworthy features). For greedy and random, the prediction is mistrusted if any untrustworthy features are present in the explanation, since these methods do not provide a notion of the contribution of each feature to the prediction. Thus for each test set prediction, we can evaluate whether the simulated user trusts it using each explanation method, and compare it to the trustworthiness oracle. <br />
<br />
Using this setup, we report the score, F1 (average of precision and recall) on the trustworthy predictions for each explanation method, averaged over 100 runs, in Table 1. The results indicate that LIME dominates others (all results are significant at p = 0.01) on both datasets, and for all of the black box models. The other methods either achieve a lower recall (i.e. they mistrust predictions more than they should) or lower precision (i.e. they trust too many predictions), while LIME maintains both high precision and high recall. Even though we artificially select which features are untrustworthy, these results indicate that LIME is helpful in assessing trust in individual predictions.<br />
<br />
[[File:table1.png|thumb|Table 1: Average F1 of trustworthiness for different explainers on a collection of classifiers and datasets.]][[File:choose_classifier.png|thumb|Figure 8: Choosing between two classifiers, as the number of instances shown to a simulated user is varied. Averages and standard errors from 800 runs.]]<br />
<br />
===Can I trust this model?===<br />
<br />
In the final simulated user experiment, we consider a case where user has to pick between 2 competing models with similar accuracy on validation dataset. For this purpose, we add 10 artificially “noisy” features. Specifically, on training and validation sets (80/20 split of the original training data), each artificial feature appears in 10% of the examples in one class, and 20% of the other, while on the test instances, each artificial feature appears in 10% of the examples in each class. This recreates the situation where the models use not only features that are informative in the real world, but also ones that introduce spurious correlations. We create pairs of competing classifiers by repeatedly training pairs of random forests with 30 trees until their validation accuracy is within 0.1% of each other, but their test accuracy differs by at least 5%. Thus, it is not possible to identify the better classifier (the one with higher test accuracy) from the accuracy on the validation data. <br />
<br />
The goal of this experiment is to evaluate whether a user can identify the better classifier based on the explanations of $B$ instances from the validation set. The simulated human marks the set of artificial features that appear in the $B$ explanations as untrustworthy, following which we evaluate how many total predictions in the validation set should be trusted (as in the previous section, treating only marked features as untrustworthy). Then, we select the classifier with fewer untrustworthy predictions, and compare this choice to the classifier with higher held-out test set accuracy. <br />
<br />
We present the accuracy of picking the correct classifier as $B$ varies, averaged over 800 runs, in Figure 8. We omit SP-parzen and RP-parzen from the figure since they did not produce useful explanations, performing only slightly better than random. LIME is consistently better than greedy, irrespective of the pick method. Further, combining sub-modular pick with LIME outperforms all other methods, in particular it is much better than RP-LIME when only a few examples are shown to the users. These results demonstrate that the trust assessments provided by SP-selected LIME explanations are good indicators of generalization."""<br />
<br />
==Conclusion and Future Work==<br />
<br />
This paper successfully argues for importance of explaining the predictions to make the model more trustworthy to the user. It proposes LIME, a comprehensive approach to faithfully explain predictions of any model in the simplest way desired. Detailed experiments demonstrating how explanations helped users pick models or discard models are provided. As a future consideration, authors want to repeat the experiments with decision trees being the interpretable model instead of sparse linear models. Authors also want to come up with an approach to perform the pick step for images. They also want to look into an automated way of selecting hyper-parameters used in LIME and SP-LIME such as K, N etc. and finally consider running the code on GPU to create explanations in realtime.<br />
<br />
==Remarks and Critique==<br />
<br />
This paper is by far the best paper I have read on model interpretability. The paper is full of new ideas and the authors clearly explain the approaches used with reasons as to why they picked it and also highlight both weaknesses and strong points of the approaches. Even though some of the tricks used are a little ambiguous on paper such as K-Lasso, super-pixels, but going through the code on GitHub, it becomes more clear. The code is also very well documented and easy to install and run to reproduce the results published. I tried running it on CIFAR-10 dataset and found it to be useful in understanding my model better.<br />
<br />
There are few minor details that are unanswered in the paper.<br />
<br />
* To interpret prediction for an image classifier, authors use super-pixels (segments) as features, but they don't give any reason as to why they picked this approach. One of the reasons I can think of for this to make sense is that convolutional neural networks learn shapes from layer to layer that become more complete moving from top to lower layers,. So in a way you can say, the classifier classifies on the basis of set of groups of pixels close to each other (as in the case of segments) than just the pixels that are spread apart far away from each other.<br />
<br />
* The authors don't say how locally faithful the interpretable model should be to the classifier with constraints on maximum number of features to be used, i.e, what is the sort of error (mean squared error in case of linear interpretable models) we are content with?<br />
<br />
* It seems that LIME can easily extended to regression cases and is not just limited to classification tasks, but the paper doesn't discuss anything in this regard.<br />
<br />
This paper is also comparable to the DeepLIFT method, which provides a global numerical recovery of sample manifold. Both methods rely heavily on the reference point chosen. LIME has the advantage of matching the behaviour of the model with interpretable explanations, while DeepLIFT focuses on the model structure and yields overall assessment of the model by manually feeding several typical reference inputs into the method.<br />
<br />
==References==<br />
[1] Marco Tulio Ribeiro, Sameer Singh, Carlos Guestrin, Model-Agnostic Interpretability of Machine Learning, presented at 2016 ICML Workshop on Human Interpretability in Machine Learning (WHI 2016), New York, NY<br />
<br />
[2] Marco Tulio Ribeiro, Sameer Singh, Carlos Guestrin, “Why Should I Trust You?” Explaining the Predictions of Any Classifier, KDD 2016 San Francisco, CA, USA<br />
<br />
[3] F. Wang and C. Rudin. Falling rule lists. In Artificial Intelligence and Statistics (AISTATS), 2015.<br />
<br />
[4] B. Efron, T. Hastie, I. Johnstone, and R. Tibshirani. Least angle regression. Annals of Statistics, 32:407–499, 2004.<br />
<br />
[5] http://scikit-image.org/docs/dev/api/skimage.segmentation.html#skimage.segmentation.quickshift<br />
<br />
[6] http://scikit-learn.org/stable/modules/generated/sklearn.linear_model.lars_path.html<br />
<br />
[7] C. Szegedy, W. Liu, Y. Jia, P. Sermanet, S. Reed, D. Anguelov, D. Erhan, V. Vanhoucke, and A. Rabinovich. Going deeper with convolutions. In Computer Vision and Pattern Recognition (CVPR), 2015.<br />
<br />
[8] A. Krause and D. Golovin. Submodular function maximization. In Tractability: Practical Approaches to Hard Problems. Cambridge University Press, February 2014.<br />
<br />
[9] J. Blitzer, M. Dredze, and F. Pereira. Biographies, bollywood, boom-boxes and blenders: Domain adaptation for sentiment classification. In Association for Computational Linguistics (ACL), 2007.<br />
<br />
[10] D. Baehrens, T. Schroeter, S. Harmeling, M. Kawanabe, K. Hansen, and K.-R. Muller. How to explain individual classification decisions. Journal of Machine Learning Research, 11, 2010.<br />
<br />
[11] https://en.wikipedia.org/wiki/Precision_and_recall<br />
<br />
[12] A. Bansal, A. Farhadi, and D. Parikh. Towards transparent systems: Semantic characterization of failure modes. In European Conference on Computer Vision (ECCV), 2014.<br />
<br />
[13] M. T. Dzindolet, S. A. Peterson, R. A. Pomranky, L. G. Pierce, and H. P. Beck. The role of trust in automation reliance. Int. J. Hum.-Comput. Stud., 58(6), 2003.<br />
<br />
[14] P. Zhang, J. Wang, A. Farhadi, M. Hebert, and D. Parikh. Predicting failures of vision systems. In Computer Vision and Pattern Recognition (CVPR), 2014.<br />
<br />
[15] Parzen Density Windows :- https://www.cs.utah.edu/~suyash/Dissertation_html/node11.html</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/_Teaching_Machines_to_Describe_Images_via_Natural_Language_Feedback&diff=30437STAT946F17/ Teaching Machines to Describe Images via Natural Language Feedback2017-11-15T14:00:06Z<p>Nkadiyar: /* Related Works */</p>
<hr />
<div>= Introduction = <br />
In the era of Artificial Intelligence, one should ideally be able to educate the robot about its mistakes, possibly without needing to dig into the underlying software. Reinforcement learning (RL) has become a standard way of training artificial agents that interact with an environment. Several works explored the idea of incorporating humans in the learning process, in order to help the reinforcement learning agent to learn faster. In most cases, the guidance comes in the form of a simple numerical (or “good”/“bad”) reward. In this work, natural language is used as a way to guide an RL agent. The author argues that a sentence provides a much stronger learning signal than a numeric reward in that we can easily point to where the mistakes occur and suggest how to correct them. <br />
<br />
Here the goal is to allow a non-expert human teacher to give feedback to an RL agent in the form of natural language, just as one would to a learning child. The author has focused on the problem of image captioning, a task where the content of an image is described using sentences. This can also be seen as a multimodal problem where the whole network/model needs to combine the solution space of learning in both the image processing and text-generation domain. Image captioning is an application where the quality of the output can easily be judged by non-experts.<br />
<br />
= Related Works =<br />
Several works incorporate human feedback to help an RL agent learn faster.<br />
#Thomaz et al. (2006) exploits humans in the loop to teach an agent to cook in a virtual kitchen. The users watch the agent learn and may intervene at any time to give a scalar reward. Reward shaping (Ng et al., 1999) is used to incorporate this information in the Markov Decision Process (MDP).<br />
#Judah et al. (2010) iterates between “practice”, during which the agent interacts with the real environment, and a critique session where a human labels any subset of the chosen actions as good or bad.<br />
#Griffith et al. (2013) proposes policy shaping which incorporates right/wrong feedback by utilizing it as direct policy labels. <br />
#Mao et. al. propose a multimodal Recurrent Neural Network (m-RNN) for image captioning on 4 crucial datasets: IAPR TC-12, Flickr 8K, Flickr 30K and MS COCO [14].Their approach involves a double network comprising of a deep RNN for sentence generation and a deep CNN for image learning.<br />
<br />
Above approaches mostly assume that humans provide a numeric reward, unlike in this work where feedback is given in natural language. A few attempts have been made to advise an RL agent using language.<br />
# Maclin et al. (1994) translated advice to a short program which was then implemented as a neural network. The units in this network represent Boolean concepts, which recognize whether the observed state satisfies the constraints given by the program. In such a case, the advice network will encourage the policy to take the suggested action.<br />
# Weston et al. (2016) incorporates human feedback to improve a text-based question answering agent.<br />
# Kaplan et al. (2017) exploits textual advice to improve training time of the A3C algorithm in playing an Atari game.<br />
<br />
The authors propose the Phrase-based Image Captioning Model which is similar to most image captioning models except that it exploits attention and linguistic information. Several recent approaches trained the captioning model with policy gradients in order to directly optimize for the<br />
desired performance metrics. This work follows the same line. <br />
<br />
There is also similar efforts on dialogue based visual representation learning and conversation modeling. These models aim to mimic human-to-human conversations while in this work the human converses with and guides an artificial learning agent.<br />
<br />
= Methodology =<br />
The framework consists of a new phrase-based captioning model trained with Policy Gradients that incorporates natural language feedback provided by a human teacher. The phrase-based captioning model allows natural guidance by a non-expert.<br />
=== Phrase-based Image Captioning ===<br />
The captioning model uses a hierarchical recurrent neural network (RNN). The model is composed of a two-level LSTM, a phrase RNN at the top level, and a word RNN that generates a sequence of words for each phrase. One can think of the phrase RNN as providing a “topic” at each time step, which instructs the word RNN what to talk about. The structure of the model is explained through Figure 1.<br />
<br />
[[File:modelham.png|center|500px|thumb|Figure 1: Hierarchical phrase-based captioning model, composed of a phrase-RNN at the top level, and a word level RNN which outputs a sequence of words for each phrase.]]<br />
<br />
<br />
A convolutional neural network is used in order to extract a set of feature vectors $a = (a_1, \dots, a_n)$, with $a_j$ a feature in location j in the input image. These feature vectors are given to the attention layer. There are also two more inputs to the attention layer, current hidden state of the phrase-RNN and output of the label unit. The label unit predicts one out of four possible phrase labels, i.e., a noun (NP), preposition (PP), verb (VP), and conjunction phrase (CP). This information could be useful for the attention layer. For example, when we have a NP the model may look at objects in the image, while for VP it may focus on more global information. Computations can be expressed with the following equations:<br />
<br />
$$<br />
\begin{align*}<br />
\small{\textrm{hidden state of the phrase-RNN at time step t}} \leftarrow h_t &= f_{phrase}(h_{t-1}, l_{t-1}, c_{t-1}, e_{t-1}) \\<br />
\small{\text{output of the label unit}} \leftarrow l_t &= softmax(f_{phrase-label}(h_t)) \\<br />
\small{\text{output of the attention layer}} \leftarrow c_t &= f_{att}(h_t, l_t, a)<br />
\end{align*}<br />
$$<br />
<br />
After deciding about phrases, the outputs of phrase-RNN go to another LSTM to produce words for each phrase. $w_{t,i}$ denotes the i-th word output of the word-RNN in the t-th phrase. There is an additional <EOP> token in word-RNN’s vocabulary, which signals the end-of-phrase. Furthermore, $h_{t,i}$ denotes the i-th hidden state of the word-RNN for the t-th phrase. <br />
$$<br />
h_{t,i} = f_{word}(h_{t,i-1}, c_t, w_{t,i}) \\<br />
w_{t,i} = f_{out}(h_{t,i}, c_t, w_{t,i-1}) \\ <br />
e_t = f_{word-phrase}(w_{t,1}, \dots ,w_{t,n})<br />
$$<br />
<br />
Note that $e_t$ encodes the generated phrase via simple mean-pooling over the words, which provides additional word-level context to the next phrase.<br />
<br />
=== Crowd-sourcing Human Feedback ===<br />
The authors have created a web interface that allows to collect feedback information. Figure 2 depicts the interface and an example of caption correction. There are two rounds of annotation. In the first round, the annotator is shown a captioned image and is asked to assess the quality of the caption, by choosing between: perfect, acceptable, grammar mistakes only, minor or major errors. They ask the annotators to choose minor and major error if the caption contained errors in semantics. They advise them to choose minor for small errors such as wrong or missing attributes or awkward prepositions, and go with major errors whenever any object or action naming is wrong. <br />
<br />
[[File:crowd.png|600px|center|thumb|Figure 2: An example of a generated caption and its corresponding feedback]]<br />
<br />
For the next (more detailed, and thus more costly) round of annotation, They only select captions which are not marked as either perfect or acceptable in the first round. Since these captions contain errors, the new annotator is required to provide detailed feedback about the mistakes. Annotators are asked to:<br />
#Choose the type of required correction (something “ should be replaced”, “is missing”, or “should be deleted”) <br />
#Write feedback in natural language (annotators are asked to describe a single mistake at a time)<br />
#Mark the type of mistake (whether the mistake corresponds to an error in object, action, attribute, preposition, counting, or grammar)<br />
#Highlight the word/phrase that contains the mistake <br />
#Correct the chosen word/phrase<br />
#Evaluate the quality of the caption after correction (it could be bad even after one round of correction)<br />
<br />
Figure 3 shows the statics of the evaluations before and after one round of correction task. The authors acknowledge the costliness of the second round of annotation.<br />
<br />
[[File:ham1.png|660px|center|thumb|Figure 3: Caption quality evaluation by the human annotators. Plot on the left shows evaluation for captions generated with the reference model (MLE). The right plot shows evaluation of the human-corrected captions (after completing at least one round of feedback).]]<br />
<br />
=== Feedback Network ===<br />
<br />
The collected feedback provides strong supervisory signal which can be used in the RL framework. In particular, the authors design a neural network (feedback network or FBN) which will provide additional reward based on the feedback sentence.<br />
<br />
RL training will require us to generate samples (captions) from the model. Thus, during training, the sampled captions for each training image will differ from the reference maximum likelihood estimation (MLE) caption for which the feedback is provided. The goal of the feedback network is to read a newly sampled caption, and judge the correctness of each phrase conditioned on the feedback. This network performs the following computations:<br />
<br />
[[File:fbn.JPG|550px|right|thumb|Figure 4: The architecture of the feedback network (FBN) that classifies each phrase in a sampled sentence (top left) as either correct, wrong or not relevant, by conditioning on the feedback sentence.]]<br />
<br />
<br />
$$<br />
h_t^{caption} = f_{sent}(h_{t-1}^{caption}, \omega_t^c) \\<br />
h_t^{feedback} = f_{sent}(h_{t-1}^{feedback}, \omega_t^f) \\<br />
q_i = f_{phrase}(\omega_{i,1}^c, \omega_{i,2}^c, \dots, \omega_{i,N}^c) \\ <br />
o_i = f_{fbn}(h_T^{caption}, h_T^{feedback }, q_i, m) \\<br />
$$<br />
<br />
<br />
Here, $\omega_t^c$ and $\omega_t^f$ denote the one-hot encoding of words in the sampled caption and feedback sentence for the t-th phrase, respectively. FBN encodes both the caption and feedback using an LSTM ($f_{sent}$), performs mean pooling ($f_{phrase}$) over the words in a phrase to represent the phrase i with $q_i$, and passes this information through a 3-layer MLP ($f_{fbn}$). The MLP accepts additional information about the mistake type (e.g., wrong object/action) encoded as a one hot vector m.<br />
<br />
=== Policy Gradient Optimization using Natural Language Feedback ===<br />
<br />
One can think of a caption decoder as an agent following a parameterized policy $p_\theta$ that selects an action at each time step. An “action” in our case requires choosing a word from the vocabulary (for the word RNN), or a phrase label (for the phrase RNN). The objective for learning the parameters of the model is the expected reward received when completing the caption $w^s = (w^s_1, \dots ,w^s_T)$. Here, $w_t^s$ is the word sampled from the model at time step t.<br />
<br />
$$<br />
L(\theta) = -\mathop{{}\mathbb{E}}_{\omega^s \sim p_\theta}[r(w^s)] <br />
$$<br />
Such an objective function is non-differentiable. Thus policy gradients are used as in [13] to find the gradient of the objective function:<br />
$$<br />
\nabla_\theta L(\theta) = - \mathop{{}\mathbb{E}}_{\omega^s \sim p_\theta}[r(w^s)\nabla_\theta \log p_\theta(w^s)]<br />
$$<br />
Which is estimated using a single Monte-Carlo sample:<br />
$$<br />
\nabla_\theta L(\theta) \approx - r(w^s)\nabla_\theta \log p_\theta(w^s)<br />
$$<br />
Then a baseline $b = r(\hat \omega)$ is used. A baseline does not change the expected gradient but can drastically reduce the variance.<br />
$$<br />
\hat{\omega}_t = argmax \ p(\omega_t|h_t) \\<br />
\nabla_\theta L(\theta) \approx - (r(\omega^s) - r(\hat{\omega}))\nabla_\theta \log p_\theta(\omega^s)<br />
$$<br />
'''Reward:''' A sentence reward is defined as a weighted sum of the BLEU scores. BLEU is an algorithm for evaluating the quality of text which has been machine-translated from one natural language to another. Quality is considered to be the correspondence between a machine's output and that of a human: "the closer a machine translation is to a professional human translation, the better it is" – this is the central idea behind BLEU. Additionally, it was one of the first metrics to claim a high correlation with human judgements of quality [10, 11 and 12] and remains one of the most popular automated and inexpensive metrics (more information about BLUE score [http://www1.cs.columbia.edu/nlp/sgd/bleu.pdf here] and a nice discussion on it [https://www.youtube.com/watch?v=ORHVgR-DVGg here]).<br />
<br />
$$<br />
r(\omega^s) = \beta \sum_i \lambda_i \cdot BLEU_i(\omega^s, ref)<br />
$$<br />
<br />
As reference captions to compute the reward, the authors either use the reference captions generated by a snapshot of the model which were evaluated as not having minor and major errors, or ground-truth captions. In addition, they weigh the reward by the caption quality as provided by the annotators (e.g. $\beta = 1$ for perfect and $\beta = 0.8$ for acceptable). They further incorporate the reward provided by the feedback network:<br />
$$<br />
r(\omega_t^p) = r(\omega^s) + \lambda_f f_{fbn}(\omega^s, feedback, \omega_t^p)<br />
$$<br />
Where $\omega^p_t$ denotes the sequence of words in the t-th phrase. Note that FBN produces a classification of each phrase. This can be converted into reward, by assigning<br />
correct to 1, wrong to -1, and 0 to not relevant. So the final gradient takes the following form:<br />
$$<br />
\nabla_\theta L(\theta) = - \sum_{p=1}^{P}(r(\omega^p) - r(\hat{\omega}^p))\nabla_\theta \log p_\theta(\omega^p)<br />
$$<br />
<br />
= Experimental Results =<br />
The authors used MS-COCO dataset. COCO is a large-scale object detection, segmentation, and captioning dataset. COCO has several features: Object segmentation, Recognition in context, Superpixel stuff segmentation, 330K images (>200K labeled), 1.5 million object instances, 80 object categories, 91 stuff categories, 5 captions per image, 250,000 people with key points. They used 82K images for training, 2K for validation, and 4K for testing. To collect feedback, they randomly chose 7K images from the training set, as well as all 2K images from validation. In addition, they use a word vocabulary size of 23,115.<br />
<br />
=== Phrase-based captioning model ===<br />
The authors analyze different instantiations of their phrase-based captioning in the following table. To sanity check, they compare it to a flat approach (word-RNN only). Overall, their model performs slightly worse (0.66 points). However, the main strength of their model is that it allows a more natural integration with feedback.<br />
<br />
[[File:table2.JPG|center]]<br />
<br />
=== Feedback network ===<br />
The authors use 9000 images to collect feedback; 5150 of them are evaluated as containing errors. Finally, they use 4174 images for the second round of annotation. They randomly select 9/10 of them to serve as a training set for feedback network, and 1/10 of them to be test set. The model achieves the highest accuracy of 74.66% when they provide it with the kind of mistake the reference caption had (e.g. an object, action, etc). This is not particularly surprising as it requires the most additional information to train the model and the most time to compile the dataset for.<br />
<br />
=== RL with Natural Language Feedback ===<br />
The following table reports the performance of several instantiations of the RL models. All models have been pre-trained using cross-entropy loss (MLE) on the full MS-COCO training set. For the next rounds of training, all the models are trained only on the 9K images.<br />
<br />
The authors define “C” captions as all captions that were corrected by the annotators and were not evaluated as containing minor or major error, and ground-truth captions for the rest of the images. For “A”, they use all captions (including captions which were evaluated as correct) that did not have minor or major errors, and GT for the rest. A detailed break-down of these captions is reported in in the following table. The authors test their model in two separate cases:<br />
<br />
*They first test a model using the standard cross-entropy loss, but which now also has access to the corrected captions in addition to the 5GT captions. This model (MLEC) is able to improve over the original MLE model by 1.4 points. They then test the RL model by optimizing the metric wrt the 5GT captions. This brings an additional point, achieving 2.4 over the MLE model. Next, the RL agent is given access to 3GT captions, the “C" captions and feedback sentences. They show that this model outperforms the no-feedback baseline by 0.5 points. If the RL agent has access to 4GT captions and feedback descriptions, a total of 1.1 points over the baseline RL model and 3.5 over the MLE model will be achieved. <br />
<br />
*They also test a more realistic scenario, in which the models have access to either a single GT caption, “C" (or “A”), and feedback. This mimics a scenario in which the human teacher observes the agent and either gives feedback about the agent’s mistakes, or, if the agent’s caption is completely wrong, the teacher writes a new caption. Interestingly, RL when provided with the corrected captions performs better than when given GT captions. Overall, their model outperforms the base RL (no feedback) by 1.2 points.<br />
<br />
[[File:table3.PNG|center]]<br />
<br />
These experiments make an important point. Instead of giving the RL agent a completely new target (caption), a better strategy is to “teach” the agent about the mistakes it is doing and suggest a correction. This is not very difficult to understand intuitively - informing the agent of its error indeed conveys more information than teaching it a completely correct answer, because the latter forces the network to "train" its memory from a sample which is, at least seemingly, insulated from its prior memory.<br />
<br />
= Conclusion =<br />
In this paper, a human teacher is enabled to provide feedback to the learning agent in the form of natural language. The authors focused on the problem of image captioning. They proposed a hierarchical phrase-based RNN as the captioning model, which allowed natural integration with human feedback.<br />
They also crowd-sourced feedback, and showed how to incorporate it in policy gradient optimization.<br />
<br />
= References=<br />
[1] Huan Ling and Sanja Fidler. Teaching Machines to Describe Images via Natural Language Feedback. In arXiv:1706.00130, 2017.<br />
<br />
[2] Shane Griffith, Kaushik Subramanian, Jonathan Scholz, Charles L. Isbell, and Andrea Lockerd Thomaz. Policy shaping: Integrating human feedback with reinforcement learning. In NIPS, 2013. <br />
<br />
[3] K. Judah, S. Roy, A. Fern, and T. Dietterich. Reinforcement learning via practice and critique advice. In AAAI, 2010. <br />
<br />
[4] A. Thomaz and C. Breazeal. Reinforcement learning with human teachers: Evidence of feedback and guidance. In AAAI, 2006. <br />
<br />
[5] Richard Maclin and Jude W. Shavlik. Incorporating advice into agents that learn from reinforcements. In National Conference on Artificial Intelligence, pages 694–699, 1994. <br />
<br />
[6] Jason Weston. Dialog-based language learning. In arXiv:1604.06045, 2016. <br />
<br />
[7] Russell Kaplan, Christopher Sauer, and Alexander Sosa. Beating atari with natural language guided reinforcement learning. In arXiv:1704.05539, 2017.<br />
<br />
[8] Andrew Y. Ng, Daishi Harada, and Stuart J. Russell. Policy invariance under reward transformations: Theory and application to reward shaping. In ICML, pages 278–287, 1999.<br />
<br />
[9] COCO Dataset http://cocodataset.org/#home<br />
<br />
[10] https://en.wikipedia.org/wiki/BLEU<br />
<br />
[11] Papineni, K., Roukos, S., Ward, T., Henderson, J and Reeder, F. (2002). “Corpus-based Comprehensive and Diagnostic MT Evaluation: Initial Arabic, Chinese, French, and Spanish Results” in Proceedings of Human Language Technology 2002, San Diego, pp. 132–137<br />
<br />
[12] Callison-Burch, C., Osborne, M. and Koehn, P. (2006) "Re-evaluating the Role of BLEU in Machine Translation Research" in 11th Conference of the European Chapter of the Association for Computational Linguistics: EACL 2006 pp. 249–256<br />
<br />
[13] Siqi Liu, Zhenhai Zhu, Ning Ye, Sergio Guadarrama, Kevin Murphy. "Improved Image Captioning via Policy Gradient optimization of SPIDEr". Under review for ICCV 2017.<br />
<br />
[14] Junhua Mao, Wei Xu, Yi Yang, Jiang Wang, Zhiheng Huang, Alan Yuille. "Deep Captioning with Multimodal Recurrent Neural Networks (m-RNN)".arXiv:1412.6632</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/_Teaching_Machines_to_Describe_Images_via_Natural_Language_Feedback&diff=30436STAT946F17/ Teaching Machines to Describe Images via Natural Language Feedback2017-11-15T13:57:29Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>= Introduction = <br />
In the era of Artificial Intelligence, one should ideally be able to educate the robot about its mistakes, possibly without needing to dig into the underlying software. Reinforcement learning (RL) has become a standard way of training artificial agents that interact with an environment. Several works explored the idea of incorporating humans in the learning process, in order to help the reinforcement learning agent to learn faster. In most cases, the guidance comes in the form of a simple numerical (or “good”/“bad”) reward. In this work, natural language is used as a way to guide an RL agent. The author argues that a sentence provides a much stronger learning signal than a numeric reward in that we can easily point to where the mistakes occur and suggest how to correct them. <br />
<br />
Here the goal is to allow a non-expert human teacher to give feedback to an RL agent in the form of natural language, just as one would to a learning child. The author has focused on the problem of image captioning, a task where the content of an image is described using sentences. This can also be seen as a multimodal problem where the whole network/model needs to combine the solution space of learning in both the image processing and text-generation domain. Image captioning is an application where the quality of the output can easily be judged by non-experts.<br />
<br />
= Related Works =<br />
Several works incorporate human feedback to help an RL agent learn faster.<br />
#Thomaz et al. (2006) exploits humans in the loop to teach an agent to cook in a virtual kitchen. The users watch the agent learn and may intervene at any time to give a scalar reward. Reward shaping (Ng et al., 1999) is used to incorporate this information in the Markov Decision Process (MDP).<br />
#Judah et al. (2010) iterates between “practice”, during which the agent interacts with the real environment, and a critique session where a human labels any subset of the chosen actions as good or bad.<br />
#Griffith et al. (2013) proposes policy shaping which incorporates right/wrong feedback by utilizing it as direct policy labels. <br />
#Mao et. al. propose a multimodal Recurrent Neural Network (m-RNN) for image captioning on 4 crucial datasets: IAPR TC-12, Flickr 8K, Flickr 30K and MS COCO [14].Their approach involves a double network comprising of a deep RNN for sentence generation and a deep CNN for image learning.<br />
<br />
Above approaches mostly assume that humans provide a numeric reward. A few attempts have been made to advise an RL agent using language.<br />
# Maclin et al. (1994) translated advice to a short program which was then implemented as a neural network. The units in this network represent Boolean concepts, which recognize whether the observed state satisfies the constraints given by the program. In such a case, the advice network will encourage the policy to take the suggested action.<br />
# Weston et al. (2016) incorporates human feedback to improve a text-based question answering agent.<br />
# Kaplan et al. (2017) exploits textual advice to improve training time of the A3C algorithm in playing an Atari game.<br />
<br />
The authors propose the Phrase-based Image Captioning Model which is similar to most image captioning models except that it exploits attention and linguistic information. Several recent approaches trained the captioning model with policy gradients in order to directly optimize for the<br />
desired performance metrics. This work follows the same line. <br />
<br />
There is also similar efforts on dialogue based visual representation learning and conversation modeling. These models aim to mimic human-to-human conversations while in this work the human converses with and guides an artificial learning agent.<br />
<br />
= Methodology =<br />
The framework consists of a new phrase-based captioning model trained with Policy Gradients that incorporates natural language feedback provided by a human teacher. The phrase-based captioning model allows natural guidance by a non-expert.<br />
=== Phrase-based Image Captioning ===<br />
The captioning model uses a hierarchical recurrent neural network (RNN). The model is composed of a two-level LSTM, a phrase RNN at the top level, and a word RNN that generates a sequence of words for each phrase. One can think of the phrase RNN as providing a “topic” at each time step, which instructs the word RNN what to talk about. The structure of the model is explained through Figure 1.<br />
<br />
[[File:modelham.png|center|500px|thumb|Figure 1: Hierarchical phrase-based captioning model, composed of a phrase-RNN at the top level, and a word level RNN which outputs a sequence of words for each phrase.]]<br />
<br />
<br />
A convolutional neural network is used in order to extract a set of feature vectors $a = (a_1, \dots, a_n)$, with $a_j$ a feature in location j in the input image. These feature vectors are given to the attention layer. There are also two more inputs to the attention layer, current hidden state of the phrase-RNN and output of the label unit. The label unit predicts one out of four possible phrase labels, i.e., a noun (NP), preposition (PP), verb (VP), and conjunction phrase (CP). This information could be useful for the attention layer. For example, when we have a NP the model may look at objects in the image, while for VP it may focus on more global information. Computations can be expressed with the following equations:<br />
<br />
$$<br />
\begin{align*}<br />
\small{\textrm{hidden state of the phrase-RNN at time step t}} \leftarrow h_t &= f_{phrase}(h_{t-1}, l_{t-1}, c_{t-1}, e_{t-1}) \\<br />
\small{\text{output of the label unit}} \leftarrow l_t &= softmax(f_{phrase-label}(h_t)) \\<br />
\small{\text{output of the attention layer}} \leftarrow c_t &= f_{att}(h_t, l_t, a)<br />
\end{align*}<br />
$$<br />
<br />
After deciding about phrases, the outputs of phrase-RNN go to another LSTM to produce words for each phrase. $w_{t,i}$ denotes the i-th word output of the word-RNN in the t-th phrase. There is an additional <EOP> token in word-RNN’s vocabulary, which signals the end-of-phrase. Furthermore, $h_{t,i}$ denotes the i-th hidden state of the word-RNN for the t-th phrase. <br />
$$<br />
h_{t,i} = f_{word}(h_{t,i-1}, c_t, w_{t,i}) \\<br />
w_{t,i} = f_{out}(h_{t,i}, c_t, w_{t,i-1}) \\ <br />
e_t = f_{word-phrase}(w_{t,1}, \dots ,w_{t,n})<br />
$$<br />
<br />
Note that $e_t$ encodes the generated phrase via simple mean-pooling over the words, which provides additional word-level context to the next phrase.<br />
<br />
=== Crowd-sourcing Human Feedback ===<br />
The authors have created a web interface that allows to collect feedback information. Figure 2 depicts the interface and an example of caption correction. There are two rounds of annotation. In the first round, the annotator is shown a captioned image and is asked to assess the quality of the caption, by choosing between: perfect, acceptable, grammar mistakes only, minor or major errors. They ask the annotators to choose minor and major error if the caption contained errors in semantics. They advise them to choose minor for small errors such as wrong or missing attributes or awkward prepositions, and go with major errors whenever any object or action naming is wrong. <br />
<br />
[[File:crowd.png|600px|center|thumb|Figure 2: An example of a generated caption and its corresponding feedback]]<br />
<br />
For the next (more detailed, and thus more costly) round of annotation, They only select captions which are not marked as either perfect or acceptable in the first round. Since these captions contain errors, the new annotator is required to provide detailed feedback about the mistakes. Annotators are asked to:<br />
#Choose the type of required correction (something “ should be replaced”, “is missing”, or “should be deleted”) <br />
#Write feedback in natural language (annotators are asked to describe a single mistake at a time)<br />
#Mark the type of mistake (whether the mistake corresponds to an error in object, action, attribute, preposition, counting, or grammar)<br />
#Highlight the word/phrase that contains the mistake <br />
#Correct the chosen word/phrase<br />
#Evaluate the quality of the caption after correction (it could be bad even after one round of correction)<br />
<br />
Figure 3 shows the statics of the evaluations before and after one round of correction task. The authors acknowledge the costliness of the second round of annotation.<br />
<br />
[[File:ham1.png|660px|center|thumb|Figure 3: Caption quality evaluation by the human annotators. Plot on the left shows evaluation for captions generated with the reference model (MLE). The right plot shows evaluation of the human-corrected captions (after completing at least one round of feedback).]]<br />
<br />
=== Feedback Network ===<br />
<br />
The collected feedback provides strong supervisory signal which can be used in the RL framework. In particular, the authors design a neural network (feedback network or FBN) which will provide additional reward based on the feedback sentence.<br />
<br />
RL training will require us to generate samples (captions) from the model. Thus, during training, the sampled captions for each training image will differ from the reference maximum likelihood estimation (MLE) caption for which the feedback is provided. The goal of the feedback network is to read a newly sampled caption, and judge the correctness of each phrase conditioned on the feedback. This network performs the following computations:<br />
<br />
[[File:fbn.JPG|550px|right|thumb|Figure 4: The architecture of the feedback network (FBN) that classifies each phrase in a sampled sentence (top left) as either correct, wrong or not relevant, by conditioning on the feedback sentence.]]<br />
<br />
<br />
$$<br />
h_t^{caption} = f_{sent}(h_{t-1}^{caption}, \omega_t^c) \\<br />
h_t^{feedback} = f_{sent}(h_{t-1}^{feedback}, \omega_t^f) \\<br />
q_i = f_{phrase}(\omega_{i,1}^c, \omega_{i,2}^c, \dots, \omega_{i,N}^c) \\ <br />
o_i = f_{fbn}(h_T^{caption}, h_T^{feedback }, q_i, m) \\<br />
$$<br />
<br />
<br />
Here, $\omega_t^c$ and $\omega_t^f$ denote the one-hot encoding of words in the sampled caption and feedback sentence for the t-th phrase, respectively. FBN encodes both the caption and feedback using an LSTM ($f_{sent}$), performs mean pooling ($f_{phrase}$) over the words in a phrase to represent the phrase i with $q_i$, and passes this information through a 3-layer MLP ($f_{fbn}$). The MLP accepts additional information about the mistake type (e.g., wrong object/action) encoded as a one hot vector m.<br />
<br />
=== Policy Gradient Optimization using Natural Language Feedback ===<br />
<br />
One can think of a caption decoder as an agent following a parameterized policy $p_\theta$ that selects an action at each time step. An “action” in our case requires choosing a word from the vocabulary (for the word RNN), or a phrase label (for the phrase RNN). The objective for learning the parameters of the model is the expected reward received when completing the caption $w^s = (w^s_1, \dots ,w^s_T)$. Here, $w_t^s$ is the word sampled from the model at time step t.<br />
<br />
$$<br />
L(\theta) = -\mathop{{}\mathbb{E}}_{\omega^s \sim p_\theta}[r(w^s)] <br />
$$<br />
Such an objective function is non-differentiable. Thus policy gradients are used as in [13] to find the gradient of the objective function:<br />
$$<br />
\nabla_\theta L(\theta) = - \mathop{{}\mathbb{E}}_{\omega^s \sim p_\theta}[r(w^s)\nabla_\theta \log p_\theta(w^s)]<br />
$$<br />
Which is estimated using a single Monte-Carlo sample:<br />
$$<br />
\nabla_\theta L(\theta) \approx - r(w^s)\nabla_\theta \log p_\theta(w^s)<br />
$$<br />
Then a baseline $b = r(\hat \omega)$ is used. A baseline does not change the expected gradient but can drastically reduce the variance.<br />
$$<br />
\hat{\omega}_t = argmax \ p(\omega_t|h_t) \\<br />
\nabla_\theta L(\theta) \approx - (r(\omega^s) - r(\hat{\omega}))\nabla_\theta \log p_\theta(\omega^s)<br />
$$<br />
'''Reward:''' A sentence reward is defined as a weighted sum of the BLEU scores. BLEU is an algorithm for evaluating the quality of text which has been machine-translated from one natural language to another. Quality is considered to be the correspondence between a machine's output and that of a human: "the closer a machine translation is to a professional human translation, the better it is" – this is the central idea behind BLEU. Additionally, it was one of the first metrics to claim a high correlation with human judgements of quality [10, 11 and 12] and remains one of the most popular automated and inexpensive metrics (more information about BLUE score [http://www1.cs.columbia.edu/nlp/sgd/bleu.pdf here] and a nice discussion on it [https://www.youtube.com/watch?v=ORHVgR-DVGg here]).<br />
<br />
$$<br />
r(\omega^s) = \beta \sum_i \lambda_i \cdot BLEU_i(\omega^s, ref)<br />
$$<br />
<br />
As reference captions to compute the reward, the authors either use the reference captions generated by a snapshot of the model which were evaluated as not having minor and major errors, or ground-truth captions. In addition, they weigh the reward by the caption quality as provided by the annotators (e.g. $\beta = 1$ for perfect and $\beta = 0.8$ for acceptable). They further incorporate the reward provided by the feedback network:<br />
$$<br />
r(\omega_t^p) = r(\omega^s) + \lambda_f f_{fbn}(\omega^s, feedback, \omega_t^p)<br />
$$<br />
Where $\omega^p_t$ denotes the sequence of words in the t-th phrase. Note that FBN produces a classification of each phrase. This can be converted into reward, by assigning<br />
correct to 1, wrong to -1, and 0 to not relevant. So the final gradient takes the following form:<br />
$$<br />
\nabla_\theta L(\theta) = - \sum_{p=1}^{P}(r(\omega^p) - r(\hat{\omega}^p))\nabla_\theta \log p_\theta(\omega^p)<br />
$$<br />
<br />
= Experimental Results =<br />
The authors used MS-COCO dataset. COCO is a large-scale object detection, segmentation, and captioning dataset. COCO has several features: Object segmentation, Recognition in context, Superpixel stuff segmentation, 330K images (>200K labeled), 1.5 million object instances, 80 object categories, 91 stuff categories, 5 captions per image, 250,000 people with key points. They used 82K images for training, 2K for validation, and 4K for testing. To collect feedback, they randomly chose 7K images from the training set, as well as all 2K images from validation. In addition, they use a word vocabulary size of 23,115.<br />
<br />
=== Phrase-based captioning model ===<br />
The authors analyze different instantiations of their phrase-based captioning in the following table. To sanity check, they compare it to a flat approach (word-RNN only). Overall, their model performs slightly worse (0.66 points). However, the main strength of their model is that it allows a more natural integration with feedback.<br />
<br />
[[File:table2.JPG|center]]<br />
<br />
=== Feedback network ===<br />
The authors use 9000 images to collect feedback; 5150 of them are evaluated as containing errors. Finally, they use 4174 images for the second round of annotation. They randomly select 9/10 of them to serve as a training set for feedback network, and 1/10 of them to be test set. The model achieves the highest accuracy of 74.66% when they provide it with the kind of mistake the reference caption had (e.g. an object, action, etc). This is not particularly surprising as it requires the most additional information to train the model and the most time to compile the dataset for.<br />
<br />
=== RL with Natural Language Feedback ===<br />
The following table reports the performance of several instantiations of the RL models. All models have been pre-trained using cross-entropy loss (MLE) on the full MS-COCO training set. For the next rounds of training, all the models are trained only on the 9K images.<br />
<br />
The authors define “C” captions as all captions that were corrected by the annotators and were not evaluated as containing minor or major error, and ground-truth captions for the rest of the images. For “A”, they use all captions (including captions which were evaluated as correct) that did not have minor or major errors, and GT for the rest. A detailed break-down of these captions is reported in in the following table. The authors test their model in two separate cases:<br />
<br />
*They first test a model using the standard cross-entropy loss, but which now also has access to the corrected captions in addition to the 5GT captions. This model (MLEC) is able to improve over the original MLE model by 1.4 points. They then test the RL model by optimizing the metric wrt the 5GT captions. This brings an additional point, achieving 2.4 over the MLE model. Next, the RL agent is given access to 3GT captions, the “C" captions and feedback sentences. They show that this model outperforms the no-feedback baseline by 0.5 points. If the RL agent has access to 4GT captions and feedback descriptions, a total of 1.1 points over the baseline RL model and 3.5 over the MLE model will be achieved. <br />
<br />
*They also test a more realistic scenario, in which the models have access to either a single GT caption, “C" (or “A”), and feedback. This mimics a scenario in which the human teacher observes the agent and either gives feedback about the agent’s mistakes, or, if the agent’s caption is completely wrong, the teacher writes a new caption. Interestingly, RL when provided with the corrected captions performs better than when given GT captions. Overall, their model outperforms the base RL (no feedback) by 1.2 points.<br />
<br />
[[File:table3.PNG|center]]<br />
<br />
These experiments make an important point. Instead of giving the RL agent a completely new target (caption), a better strategy is to “teach” the agent about the mistakes it is doing and suggest a correction. This is not very difficult to understand intuitively - informing the agent of its error indeed conveys more information than teaching it a completely correct answer, because the latter forces the network to "train" its memory from a sample which is, at least seemingly, insulated from its prior memory.<br />
<br />
= Conclusion =<br />
In this paper, a human teacher is enabled to provide feedback to the learning agent in the form of natural language. The authors focused on the problem of image captioning. They proposed a hierarchical phrase-based RNN as the captioning model, which allowed natural integration with human feedback.<br />
They also crowd-sourced feedback, and showed how to incorporate it in policy gradient optimization.<br />
<br />
= References=<br />
[1] Huan Ling and Sanja Fidler. Teaching Machines to Describe Images via Natural Language Feedback. In arXiv:1706.00130, 2017.<br />
<br />
[2] Shane Griffith, Kaushik Subramanian, Jonathan Scholz, Charles L. Isbell, and Andrea Lockerd Thomaz. Policy shaping: Integrating human feedback with reinforcement learning. In NIPS, 2013. <br />
<br />
[3] K. Judah, S. Roy, A. Fern, and T. Dietterich. Reinforcement learning via practice and critique advice. In AAAI, 2010. <br />
<br />
[4] A. Thomaz and C. Breazeal. Reinforcement learning with human teachers: Evidence of feedback and guidance. In AAAI, 2006. <br />
<br />
[5] Richard Maclin and Jude W. Shavlik. Incorporating advice into agents that learn from reinforcements. In National Conference on Artificial Intelligence, pages 694–699, 1994. <br />
<br />
[6] Jason Weston. Dialog-based language learning. In arXiv:1604.06045, 2016. <br />
<br />
[7] Russell Kaplan, Christopher Sauer, and Alexander Sosa. Beating atari with natural language guided reinforcement learning. In arXiv:1704.05539, 2017.<br />
<br />
[8] Andrew Y. Ng, Daishi Harada, and Stuart J. Russell. Policy invariance under reward transformations: Theory and application to reward shaping. In ICML, pages 278–287, 1999.<br />
<br />
[9] COCO Dataset http://cocodataset.org/#home<br />
<br />
[10] https://en.wikipedia.org/wiki/BLEU<br />
<br />
[11] Papineni, K., Roukos, S., Ward, T., Henderson, J and Reeder, F. (2002). “Corpus-based Comprehensive and Diagnostic MT Evaluation: Initial Arabic, Chinese, French, and Spanish Results” in Proceedings of Human Language Technology 2002, San Diego, pp. 132–137<br />
<br />
[12] Callison-Burch, C., Osborne, M. and Koehn, P. (2006) "Re-evaluating the Role of BLEU in Machine Translation Research" in 11th Conference of the European Chapter of the Association for Computational Linguistics: EACL 2006 pp. 249–256<br />
<br />
[13] Siqi Liu, Zhenhai Zhu, Ning Ye, Sergio Guadarrama, Kevin Murphy. "Improved Image Captioning via Policy Gradient optimization of SPIDEr". Under review for ICCV 2017.<br />
<br />
[14] Junhua Mao, Wei Xu, Yi Yang, Jiang Wang, Zhiheng Huang, Alan Yuille. "Deep Captioning with Multimodal Recurrent Neural Networks (m-RNN)".arXiv:1412.6632</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/_Coupled_GAN&diff=30154STAT946F17/ Coupled GAN2017-11-13T23:46:40Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div><br />
This is a summary of NIPS 2016 paper [1].<br />
== Introduction==<br />
<br />
Generative models attempt to characterize and estimate the underlying probability distribution of the data (typically images) and in doing so generate samples from the aforementioned learned distribution. Moment-matching generative networks, Variational auto-encoders, and Generative Adversarial Networks (GANs) are some of the most popular (and recent) class of techniques in this burgeoning literature on generative models. The authors of the paper we are reviewing focus on proposing an extension to the class of GANs.<br />
<br />
The novelty of the proposed Coupled GAN (CoGAN) method lies in extending the GAN procedure (described in the next section) to the multi-domain setting. That is, the CoGAN methodology attempts to learn the (underlying) joint probability distribution of multi-domain images as a natural extension from the marginal setting associated with the vanilla GAN framework. This is inspired by the idea that deep neural networks learn a hierarchical feature representation. Given the dense and active literature on generative models, generating images in multiple domains is far from groundbreaking. Related works revolve around multi-modal deep learning ([2],[3]), semi-coupled dictionary learning ([4]), joint embedding space learning ([5]), cross-domain image generation ([6],[7]) to name a few. Thus, the novelty of the authors' contributions to this field comes from two key differentiating points. Firstly, this was (one of) the first papers to endeavor to generate multi-domain images with the GAN framework. Secondly, and perhaps more significantly, the authors proposed to learn the underlying joint distribution without requiring the presence of tuples of corresponding images in the training set. Only sets of images drawn from the (marginal) distributions of the separate domains is sufficient. As per the authors' claim, constructing tuples of corresponding images to train from is challenging and a potential bottleneck for multi-domain image generation. One way around this bottleneck is thus to use their proposed CoGAN methodology. More details of how the authors achieve joint-distribution learning will be provided in the Coupled GAN section below.<br />
<br />
== Generative Adversarial Networks==<br />
<br />
A typical GAN framework consists of a generative model and a discriminative model. The generative model, which often is a de-convolutional network, takes as input a random ''latent'' vector (typically uniform or Gaussian) and synthesizes novel images resembling the real images (training set). The discriminative model, often a convolutional network, on the other hand, tries to distinguish between the fake synthesized images and the real images. The idea then is to let the two component models of the GAN framework "compete" with each other in the form of a min-max two player game. <br />
<br />
To further clarify and fix this idea, we introduce the mathematical setup of GANs following the notation used by the authors of this paper for sake of consistency. Let us define the following in our setup:<br />
<br />
:<math> \mathbf{x}-</math> natural image drawn from underlying distribution <math> p_X</math>,<br />
:<math> \mathbf{z} \sim U[-1,1]^d-</math> a latent random vector,<br />
:$g-$ generative model, $f-$ discriminative model.<br />
<br />
Ideally we are aiming for the system of these two ''adversarial'' networks to behave as:<br />
:Generator: $g(\mathbf{z})$ outputs an image with same support as $\mathbf{x}$. The probability density of the images output by $g$ can be denoted by $p_G$,<br />
:Discriminator: $f(\mathbf{x})=1$ if $\mathbf{x} \sim p_X$ and $f(\mathbf{x})=0$ if $\mathbf{x} \sim p_G$.<br />
<br />
To train such a system of networks given our goal (i.e., $p_G \rightarrow p_X$) we must treat such a framework as the following minimax two player game:<br />
<br />
$\displaystyle \max_{g}$<br />
$\min\limits_{f} V(g,f) = \mathop{\mathbb{E}}_{\mathbf{x} \sim p_X}[-\log(f(\mathbf{x}))] + \mathop{\mathbb{E}}_{\mathbf{z} \sim p_{Z}(\mathbf{z})}[-\log(1-f(g(\mathbf{z})))] $.<br />
<br />
See [8], the seminal paper on this topic, for more information.<br />
<br />
Some of the crucial advantages of GANs are that Markov chains are never needed, only backprop is used to obtain gradients, no inference is needed during<br />
learning, and a wide variety of functions can be incorporated into the model [16].<br />
<br />
== Coupled Generative Adversarial Networks==<br />
<br />
The overarching goal of this framework is to learn a joint distribution of multi-domain images from data. That is, a density value is assigned to each joint occurrence of images in different domains. Examples of such pair of images in different domains include images of a particular scene with different modalities (color and depth) or images of the same face but with different facial attributes. <br />
<br />
To this end, the CoGAN setup consists of a pair of GANs, denoted as $GAN_1$ and $GAN_2$. Each GAN is tasked with synthesizing images in one domain. A naive training of such a system will result in learning the product of the two marginal distributions i.e., independence. However, by forcing the two GANs to share weights, the authors were able to demonstrate that they could in ''some sense'' learn the joint distribution of images. We will now describe the details of the generator and discriminator components of the setup and conclude this section with a summary of CoGAN learning algorithm.<br />
<br />
===Generator Models===<br />
<br />
Suppose $\mathbf{x_1} \sim p_{X_1}$ and $\mathbf{x_2} \sim p_{X_2}$ denote the natural images being drawn from the two marginal distributions of <br />
domain 1 and domain 2. Further, let $g_1$ be the generator of $GAN_1$ and $g_2$ be the generator of $GAN_2$. Both these generators take as input the latent vector $\mathbf{z}$ as defined in the previous section as input and out images in their specific domains. For completeness, denote the distributions of $g_1(\mathbf{z})$ and $g_2(\mathbf{z})$ as $p_{G_1}$ and $p_{G_2}$ respectively. We can characterize these two generator models as multi-layer perceptrons in the following way:<br />
<br />
\begin{align*}<br />
g_1(\mathbf{z})=g_1^{(m_1)}(g_1^{(m_1 -1)}(\dots g_1^{(2)}(g_1^{(1)}(\mathbf{z})))), \quad g_2(\mathbf{z})=g_2^{(m_2)}(g_2^{(m_2-1)}(\dots g_2^{(2)}(g_2^{(1)}(\mathbf{z})))),<br />
\end{align*}<br />
where $g_1^{(i)}$ and $g_2^{(i)}$ are the $i^{th}$ layers of $g_1$ and $g_2$ which respectively have a total of $m_1$ and $m_2$ layers each. Note $m_1$ need not be the same as $m_2$.<br />
<br />
As the generator networks can be thought of as an inverse of the prototypical convolutional networks (just as an example), the layers of these generator networks gradually decode information from high-level abstract concepts to low-level details(last few layers). Taking this idea as the blueprint for the inner-workings of generator networks, the author's hypothesize that corresponding images in two domains share the same high-level semantics but with differing lower-level details. To put this hypothesis to practice, they forced the first $k$ layers of $g_1$ and $g_2$ to have identical structures and share the same weights. That is, $\mathbf{\theta}_{g_1^{(i)}}=\mathbf{\theta}_{g_2^{(i)}}$ for $i=1,\dots,k$ where $\mathbf{\theta}_{g_1^{(i)}}$ and $\mathbf{\theta}_{g_1^{(i)}}$ represents the parameters of the layers $g_1^{(i)}$ and $g_2^{(i)}$ respectively. Hence the two generator networks share the starting $k$ layers of the deep network and have different last layers to decode the differing material details in each domain.<br />
<br />
===Discriminative Models===<br />
<br />
Suppose $f_1$ and $f_2$ are the respective discriminative models of the two GANs. These models can be characterized by <br />
\begin{align*}<br />
f_1(\mathbf{x}_1)=f_1^{(n_1)}(f_1^{(n_1 -1)}(\dots f_1^{(2)}(f_1^{(1)}(\mathbf{x}_1)))), \quad f_2(\mathbf{x}_2)=f_2^{(n_2)}(f_2^{(n_2-1)}(\dots f_2^{(2)}(f_2^{(1)}(\mathbf{x}_1)))),<br />
\end{align*}<br />
where $f_1^{(i)}$ and $f_2^{(i)}$ are the $i^{th}$ layers of $f_1$ and $f_2$ which respectively have a total of $n_1$ and $n_2$ layers each. Note $n_1$ need not be the same as $n_2$. In contrast to generator models, the first layers of $f_1$ and $f_2$ extract the lower level details where the last layers extract the abstract higher level details. To reflect the prior hypothesis of shared higher level semantics between corresponding images, we can force $f_1$ and $f_2$ to now share the weights for last $l$ layers. That is, $\mathbf{\theta}_{f_1^{(n_1-i)}}=\mathbf{\theta}_{f_2^{(n_2-i)}}$ for $i=0,\dots,l-1$ where $\mathbf{\theta}_{f_1^{(i)}}$ and $\mathbf{\theta}_{f_1^{(i)}}$ represents the parameters of the layers $f_1^{(i)}$ and $f_2^{(i)}$ respectively. Unlike in the generative models, weight sharing in the discriminative models is not essential to estimating the joint distribution of images, however it is beneficial by reducing the total number of parameters in the network.<br />
<br />
===Coupled GAN (CoGAN) Framework and Learning===<br />
The following figure taken from the paper summarizes the system of models described in the previous subsections. <br />
<center><br />
[[File:CoGAN-1.PNG]]<br />
</center><br />
The CoGAN framework can be expressed as the following constrained min-max game<br />
<br />
\begin{align*}<br />
\max\limits_{g_1,g_2} \min\limits_{f_1, f_2} V(f_1,f_2,g_1,g_2)\quad \text{subject to} \ \mathbf{\theta}_{g_1^{(i)}}=\mathbf{\theta}_{g_2^{(i)}}, i=1,\dots k, \quad \mathbf{\theta}_{f_1^{(n_1-j)}}=\mathbf{\theta}_{f_2^{(n_2-j)}}, j=1,\dots,l-1, <br />
\end{align*}<br />
where the value function V is characterized as <br />
\begin{align*}<br />
\mathop{\mathbb{E}}_{\mathbf{x}_1 \sim p_{X_1}}[-\log(f_1(\mathbf{x_1}))] + \mathop{\mathbb{E}}_{\mathbf{z} \sim p_{Z}(\mathbf{z})}[-\log(1-f_1(g_1(\mathbf{z})))]+\mathop{\mathbb{E}}_{\mathbf{x}_2 \sim p_{X_2}}[-\log(f_2(\mathbf{\mathbf{x}_2}))] + \mathop{\mathbb{E}}_{\mathbf{z} \sim p_{Z}(\mathbf{z})}[-\log(1-f_2(g_2(\mathbf{z})))]. <br />
\end{align*}<br />
<br />
For the purposes of storytelling, we can describe this game to have two teams with two players each. The generative models are in the same team and collaborate with each other to synthesize a pair of images in two different domains with goal of fooling the discriminative models. The discriminative models than with collaboration try to differentiate between images drawn from the training data in their respective domains and the images generated by the respective generative models. The training algorithm for the CoGAN that was used is described in the following figure.<br />
<center><br />
[[File:CoGAN-2.PNG]]<br />
</center><br />
<br />
'''Important Remarks:''' <br />
<br />
CoGAN learning requires training samples drawn from the marginal distributions, $p_{X_1}$ and $p_{X_2}$ . It does not rely on samples drawn from the joint distribution, $p_{X_1,X_2}$ , where corresponding supervision would be available. Here, the main contribution is in showing that with just samples drawn separately from the marginal distributions, CoGAN can learn a joint distribution of images in the two domains. Both weight-sharing constraint and adversarial training are essential for enabling this capability. <br />
<br />
Unlike autoencoder learning which encourages the generated image pair to be identical to the target pair , the adversarial training only encourages the generated pair of images to be individually resembling the images in the respective domains, and ignore the correlation between them. Shared parameters, on the other hand, contribute to matching the correlation: the neurons responsible for decoding high-level semantics can be shared to produce highly correlated image pairs.<br />
<br />
==Experiments==<br />
<br />
To begin with, note that the authors do not use corresponding images in the training set in accordance with the goal of ''learning'' the joint distribution of multi-domain images without correspondence supervision. As at the time the paper was written, there were no existing approached with identical prerogatives (i.e., training with no correspondence supervision), they compared CoGAN with conditional GAN (see [10]) for more details on conditional GAN).<br />
<br />
===MNIST Dataset===<br />
<br />
The MNIST training set was experimented with two tasks:<br />
<br />
# Task A: Learning a joint distribution of a digit and its edge image. <br />
# Task B: Learning a joint distribution of a digit and its negative image. <br />
<br />
For the generative models, the authors used convolutional networks with 5 identical layers. They varied the number of shared layers as part of their experimental setup. The two discriminative models were a version of the LeNet ([9]). The results of the CoGAN generation scheme are displayed in the figure below.<br />
<br />
<center><br />
[[File:CoGAN-3.PNG]]<br />
</center><br />
As you can see from the figure above, the CoGAN system was able to generate pairs of corresponding images without explicitly training with correspondence supervision. This was naturally due to sharing weights in lower levels used for decoding high-level semantics. Without sharing these weights, the CoGAN would just output a pair of unrelated images in the two domains.<br />
<br />
To investigate the effects of weight sharing in the generator/discriminator models used for both tasks, the authors varied the number of shared levels. To quantify the performance of the generator, the image generated by $GAN_1$ (domain 1) was transformed to the 2nd domain using the same method used to generate training images in the 2nd domain. Then this transformed image was compared with the image generated by $GAN_2$. Naturally, if the joint distribution was learned completely, these two images would be identical. With that goal in mind, the authors used pixel agreement ratios for 10000 images as the evaluation metric. In particular, 5 trails with different weight initializations were used and an average pixel agreement ratio was taken. The results depicting the relationship between the average pixel agreement ratio and a number of shared layers are summarized in the figure below. <br />
<center><br />
[[File:CoGAN-4.PNG]]<br />
</center><br />
The results naturally offered some corroboration to our intuitions. The greater the number of shared layers in the generator models, the higher the pixel agreement ratios. Interestingly the number of shared layers in the discriminative model does not seem to affect the pixel agreement ratios. Note this is a pretty naive and toy example as we by nature of the evaluation criteria have a deterministic way of generating an image in the 2nd domain.<br />
<br />
Finally, for this example, the authors compared the CoGAN framework with the conditional GAN model. For the conditional GAN, the generative and discriminative models were identical to those used for the CoGAN results. The conditional GAN additionally took a binary variable (the conditioning variable) as input. When the binary variable was 0, the conditional GAN synthesized an image in domain 1 and when it was 1 an image in domain 2. Naturally, for a fair comparison, the training set did not already contain corresponding pairs of images. The experiments were conducted for the two tasks described above and the pixel agreement ratio (PAR) was the evaluation criteria. For Task A, the CoGAN resulted in a PAR of 0.952 in comparison with 0.909 for the conditional GAN. For Task B, the CoGAN resulted in a PAR of 0.967 compared with a PAR of 0.778 for the conditional GAN. The results are not particularly eye opening as the CoGAN was more specifically designed for purpose of learning the joint distribution of multi-domain images whereas these tasks are just a very niche application for the conditional GAN. Nevertheless, for Task B, the results look promising.<br />
<br />
=== CelebFaces Attributes Dataset===<br />
<br />
For this experiment, the authors trained the CoGAN, using the CelebFaces Attributes Dataset, to generate pairs of faces with an attribute (domain 1) and without the attribute (domain 2). CelebFaces Attributes Dataset (CelebA) is a large-scale face attributes dataset with more than 200K celebrity images, each with 40 attribute annotations. The images in this dataset cover large pose variations and background clutter. CelebA has large diversities, large quantities, and rich annotations, including 10,177 number of identities, 202,599 number of face images, and 5 landmark locations, 40 binary attributes annotations per image.<br />
Convolutional networks with 7 layers for both the generative and discriminative models were used. The dataset contains a large variety of poses and background clutter. Attributes can include blonde/non-blonde hair, smiling/ not smiling, or with/without sunglasses for example. The resulting synthesized pair of images are shown in the figure below realized as spectrum traveling from one point to another (resembles changing faces).<br />
<br />
<center><br />
[[File:CoGAN-5.PNG]]<br />
</center><br />
<br />
=== Color and Depth Images===<br />
For this experiment the authors used two sources: the RGBD dataset and the NYU dataset. The two image domains (color-domain 1 and depth-domain 2) under consideration are the same for both the datasets. As usual, no corresponding images were fed into the training of the CoGAN framework. The resulting rendering of pairs of color and depth images for both the datasets are depicted in the figure below. <br />
<br />
<center><br />
[[File:CoGAN-6.PNG]]<br />
</center><br />
<br />
As is evident from the images, through the sharing of layers during training, the CoGAN was able to learn the appearance-depth correspondence.<br />
<br />
== Applications==<br />
<br />
===Unsupervised Domain Adaptation (UDA)===<br />
<br />
UDA involves adapting a classifier trained in one domain to conduct a classification task in a new domain which only contains ''unlabeled'' training data which disqualifies re-training of the classifier in this new domain. Some prior work in the field includes subspace learning ([11],[12]) and deep discriminative network learning ([13],[14]). The authors in the paper experimented with the MNIST and USPS datasets to showcase the applicability of the CoGAN framework for the UDA problem. A similar network architecture as was used for the MNIST experiment was employed for this application. The MNIST and USPS datasets has been denoted as $D_1$ and $D_2$ respectively in the paper. In accordance with the problem specification, no label information was used from $D_2$. <br />
<br />
The CoGAN is trained by jointly solving the classification problem in the MNIST domain using the labels provided in $D_1$ and the CoGAN learning problem which uses images in both $D_1$ and $D_2$. This training process produces two classifiers. That is, $c_1(x_1)=c(f_1^{(3)}(f_1^{(2)}(f_1^{(1)}(x_1))))$ for MNIST and $c_2(x_2)=c(f_2^{(3)}(f_2^{(2)}(f_2^{(1)}(x_2))))$ and USPS. Note $c()$ here denotes the softmax layer which is added on top of the other layers in the respective discriminative networks of the GANs. Further due to weight sharing the last two layers of the discriminative model would be identical and have the same weights. The classifier $c_2$ is then used for digit classification in the USPS dataset. The author reported a 91.2% average accuracy when classifying the USPS dataset. The mirror problem of classifying the MNIST dataset without labels using the fully characterized USPS dataset achieved an average accuracy of 89.1%. These results appear to significantly outperform (prior top classification accuracy lies roughly around 60-65%) what the authors ''claim'' to be the state of the art methods in the UDA literature. In particular the state of the art was noted to be described in [20].<br />
<br />
===Cross-Domain Image Transformation===<br />
<br />
Let $\mathbf{x}_1$ be a image in the 1st domain. The goal then is to find a corresponding image $\mathbf{x}_2$ in 2nd domain such that the joint probability density $p(\mathbf{x}_1,\mathbf{x}_2)$ is maximized. Given the two generators $g_1$ and $g_2$, one can achieve the cross-domain transformation by first finding the latent random vector that generates the input image $\mathbf{x}_1$ in the 1st domain. This amounts to the optimization: $\mathbf{z}^{*}=argmin_{\mathbf{z}}L(g_1(\mathbf{z}),\mathbf{x}_1)$, where $L$ is the loss function measuring the difference/distance between the two images. After finding $z^*$, one can apply $g_2$ to obtain the transformed image, $x_2 = g_2(z^*)$. Some very preliminary results were provided by the authors but there is nothing noteworthy to warrant discussion. The figure depicting their results is provided below for sake of completeness.<br />
[[File:CoGAN-7.PNG]]<br />
<br />
===Eyewitness Facial Composite Generation===<br />
Since a Coupled GAN requires only a small set of images acquired separately from the marginal distributions of the individual domains, forensic criminal feature generation could find use for CoGAN. Due to the high noise and variation in facial composites derived from eyewitness statements (marginal), CoGAN could assist in narrowing down the search space.<br />
<br />
== Discussion and Summary==<br />
The main contributions of the paper can be summarized as:<br />
<br />
# A CoGAN framework for learning a joint distribution of multi-domain images was proposed. <br />
# The training is achieved by a simple weight sharing scheme for the generative and discriminative networks in the absence of any correspondence supervision in the training set. This can be construed as learning the joint distribution by using samples from the marginal distribution of images. <br />
# The experiments with digits, faces, and color/depth images provided some corroboration that the CoGAN system could synthesize corresponding pairs of images. <br />
# An application of the CoGAN framework for the problem of Unsupervised Domain Adaptation (UDA) was introduced. The preliminary results appear to be extremely promising for the task of adapting digit classifiers from MNIST to USPS data and vice-versa. <br />
#An application for the task of cross-domain image transformation was hypothesized with some very basic proof of concept results provided. <br />
# The setup naturally lends to more than the two domain setting focused on in the paper for experimental purposes. <br />
<br />
While the summary provided above adopted an objective filter, the following list enumerates the major ''subjective'' critical review points for this paper:<br />
# It appears the authors took components of various well-established techniques in the literature and produced the CoGAN framework. Weight-sharing is a well documented idea as was correspondance/ multi-modal learning along with the GAN problem formation and training. However, when the components are put together in this way, they form a modest and timely novel contribution to the literature of generative networks.<br />
# With such prominent preliminary results for the problem of UDA, the authors could have provided some additional details of their training procedure (slightly unclear) and additional experiments under the UDA umbrella to fortify what appears to be a ''ground breaking'' result when compared with state of the art methods. <br />
# The cross-domain image transformation application example was almost an after thought. More details could have been provided in the supplementary file if pressed for space or perhaps just merely relegated to a follow-up paper/work.<br />
<br />
== References and Supplementary Resources==<br />
:[1] Liu, Ming-Yu, and Oncel Tuzel. "Coupled generative adversarial networks." Advances in neural information processing systems. 2016.<br />
:[2] Srivastava, Nitish, and Ruslan R. Salakhutdinov. "Multimodal learning with deep boltzmann machines." Advances in neural information processing systems. 2012.<br />
:[3] Ngiam, Jiquan, et al. "Multimodal deep learning." Proceedings of the 28th international conference on machine learning (ICML-11). 2011.<br />
:[4] Wang, Shenlong, et al. "Semi-coupled dictionary learning with applications to image super-resolution and photo-sketch synthesis." Computer Vision and Pattern Recognition (CVPR), 2012 IEEE Conference on. IEEE, 2012.<br />
:[5] Kiros, Ryan, Ruslan Salakhutdinov, and Richard S. Zemel. "Unifying visual-semantic embeddings with multimodal neural language models." arXiv preprint arXiv:1411.2539 (2014).<br />
:[6] Yim, Junho, et al. "Rotating your face using multi-task deep neural network." Proceedings of the IEEE Conference on Computer Vision and Pattern Recognition. 2015.<br />
:[7] Reed, Scott E., et al. "Deep visual analogy-making." Advances in neural information processing systems. 2015. <br />
:[8] Goodfellow, Ian, et al. "Generative adversarial nets." Advances in neural information processing systems. 2014.<br />
:[9] LeCun, Yann, et al. "Gradient-based learning applied to document recognition." Proceedings of the IEEE 86.11 (1998): 2278-2324.<br />
:[10] Mirza, Mehdi, and Simon Osindero. "Conditional generative adversarial nets." arXiv preprint arXiv:1411.1784 (2014).<br />
:[11] Long, Mingsheng, et al. "Transfer feature learning with joint distribution adaptation." Proceedings of the IEEE international conference on computer vision. 2013.<br />
:[12] Fernando, Basura, Tatiana Tommasi, and Tinne Tuytelaars. "Joint cross-domain classification and subspace learning for unsupervised adaptation." Pattern Recognition Letters 65 (2015): 60-66. <br />
:[13] Tzeng, Eric, et al. "Deep domain confusion: Maximizing for domain invariance." arXiv preprint arXiv:1412.3474 (2014).<br />
:[14] Rozantsev, Artem, Mathieu Salzmann, and Pascal Fua. "Beyond sharing weights for deep domain adaptation." arXiv preprint arXiv:1603.06432 (2016).<br />
:[15] http://mmlab.ie.cuhk.edu.hk/projects/CelebA.html<br />
:[16] https://arxiv.org/pdf/1406.2661.pdf<br />
<br />
Implementation Example: https://github.com/mingyuliutw/cogan</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Imagination-Augmented_Agents_for_Deep_Reinforcement_Learning&diff=30153Imagination-Augmented Agents for Deep Reinforcement Learning2017-11-13T23:34:05Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>=Introduction=<br />
An interesting research area in reinforcement learning is developing intelligent agents for playing video games. Before the introduction of deep learning, video game agents were commonly coded based on Monte-Carlo Tree Search of pre-set rules. Recent research has shown deep reinforcement learning to be very successful at playing video games like Atari 2600. To be specific, the method (Figure 1) is called Deep Q-Learning (DQN) which learns the optimal actions based on current observations (raw pixels) [[#Reference|[Mnih et al., (2015)]]]. However, there are some complex games where DQN fails to learn: some games need to solve a sub-problem without explicit reward or contain irreversible domains, where actions can be catastrophic. A typical example of these games is [https://en.wikipedia.org/wiki/Sokoban Sokoban]. Even as humans are playing the game, planning and inference are needed. This kind of game raises challenges to RL.<br />
<br />
[[File:DQN.png|800px|center|thumb|Figure 1: Deep Q-Learning Architecture]]<br />
<br />
In Reinforcement Learning, the algorithms can be divided into two categories: '''model-free''' algorithm and '''model-based''' algorithm. The model-based reinforcement learning tries to infer environment to gain the reward while model-free reinforcement learning does not use the environment to learn the action that results in the best reward. More specifically, model-based methods learn the model (The Reward function: $R(s, s^{'})$ and the Transition probability $P(s^{'} | s, a)$ ) of the environment, while model-free methods never explicitly learn the model of the environment. DQN, mentioned above(Figure 1), is a model-free method. It takes raw pixels as input and maps them to values or actions. As a drawback, large amounts of training data is required. In addition, the policies are not generalized to new tasks in the same environment. A model-based method is trying to build a model for the environment. By querying the model, agents can avoid irreversible, poor decisions. As an approximation of the environment, it can enable better generalization across states. However, this method only shows success in limited settings, where an exact transition model is given or in simple domains. In complex environments, model-based methods suffer from model errors from function approximation. These errors compound during planning, causing poor agent performance. Currently, there is no model-based method that is robust against imperfections.<br />
<br />
In this paper, the authors introduce a novel deep reinforcement learning architecture called Imagination-Augmented Agents(I2As). Literally, this method enables agents to learn to interpret predictions from a learned environment model to construct implicit plans. It is a combination of model-free and model-based aspects. The advantage of this method is that it learns in an end-to-end way to extract information from model simulations without making any assumptions about the structure or the perfections of the environment model.<br />
As shown in the results, this method outperforms DQN in the games: Sokoban, and MniniPacman. In addition, the experiments all show that I2A is able to successfully use imperfect models.<br />
<br />
=Motivation=<br />
Although the structure of this method is complex, the motivation is intuitive: since the agent suffers from irreversible decisions, attempts in simulated states may be helpful.<br />
To improve the expensive search space in traditional MCTS methods, adding decision from policy network can reduce search steps. In order to keep context information, rollout results are encoded by an LSTM encoder. The final output is combining the result from model-free network and model-based network.<br />
<br />
=Related Work=<br />
There are some works that try to apply deep learning to model-based reinforcement learning. The popular approach is to learn a neural network from the environment and apply the network in classical planning algorithms. These works can not handle the mismatch between the learned model and the ground truth. Liu et al.(2017) use context information from trajectories, but in terms of imitation learning.<br />
<br />
To deal with imperfect models, Deisenroth and Rasmussen(2011) try to capture model uncertainty by applying high-computational Gaussian Process models. In order to develop such a policy search method, the authors of this paper used analytic gradients of an approximation to the expected return for indirect policy search. This means by learning a probabilistic dynamics model and explicitly incorporating model uncertainty into long-term planning, this policy search method can cope with very little data and facilitates learning from scratch in only a few trials. <br />
<br />
Similar ideas can be found in a study by Hamrick et al.(2017): they present a neural network that queries expert models, but just focus on meta-control for continuous contextual bandit problems. Pascanu et al.(2017) extend this work by focusing on explicit planning in sequential environments.<br />
<br />
=Approach=<br />
The summary of the architecture of I2A can be seen in Figure 2.<br />
[[File:i2a.png|800px|center|thumb|Figure 2: The Architecture of I2A]]<br />
The observation $O_t$ (Figure 2 right) is fed into two paths, the model-free path is just common DQN which predicts the best action given $O_t$, whereas the model-based path performs a rollout strategy, the aggregator combines the $n$ rollout encoded outputs($n$ equals to the number of actions in the action space), and forwards the results to next layer. Together they are used to generate a policy function $\pi$ to output an action. In each rollout operation, the imagination core is used to predict the future state and reward.<br />
<br />
===Imagination Core===<br />
The imagination-augmented agents adopt a concept called the "imagination encoder" which is a neural network which learns to extract relevant information that impacts the agent's future decisions, ignoring those information that is irrelevant. In particular, these agents have the following features: (i) they have the ability to learn to interpret their internal simulations which captures the environmental dynamics, (ii) they adapt to the number of imagined trajectors which makes the imagination more efficient, and finally (iii) they have the ability to learn different strategies to construct plans by choosing the appropriate trajectory. The imagination core(Figure 2 left) is the key role in the model-based path. It consists of two parts: environment model and rollout policy. The former is an approximation of the environment and the latter is used to simulate imagined trajectories, which are interpreted by a neural network and provided as additional context to a policy network. <br />
<br />
====environment model====<br />
In order to augment agents with imagination, the method relies on environment models that, given current information, can be queried to make predictions about the future. In this work, the environment model is built based on action-conditional next-step predictors, which receive input contains current observation and current action, and predict the next observation and the next reward(Figure 3).<br />
[[File:environment model.png|800px|center|thumb|Figure 3: Environment Model]]<br />
<br />
The authors can either pretrain the environment model before embedding it (with frozen weights) within the I2A architecture, or jointly train it with the agent by adding $l_{model}$ to the total loss as an auxiliary loss. In practice, they found that pre-training the environment model led to faster<br />
runtime of the I2A architecture, so they adopted this strategy.<br />
<br />
====rollout policy====<br />
The rollout process is regarded as the simulated trajectories. In this work, the rollout is performed for each possible action in the environment. <br />
<br />
A rollout policy $\hat \pi$ is a function that takes current observation $O$ and outputs an action $a$ that potentially leads to maximal reward. In this architecture, the rollout policy can be a DQN network. In the experiment, the rollout pocily$\hat \pi$ is broadcasted and shared. After experiments on the types of rollout policies(random, pre-trained), the authors found the efficient strategy is to distill the policy into a model-free policy, which consists in creating a smaill model-free network $\hat \pi(O_t)$, and adding to the total loss a cross entropy auxiliary loss between the imagination-augmented policy $\pi(O_t)$ as computed on the current observation, and the policy $\hat \pi(O_t)$ as computed on the same observation.<br />
<br />
$$<br />
l_{dist} (\pi, \hat \pi)(O_t) = \lambda_{dist} \sum_a \pi(a|O_t)log(\hat \pi(a|O_t))<br />
$$<br />
<br />
Together as the imagination core, these two parts produces $n$ trajectories $\hat \tau_1,...,\hat \tau_n$. Each imagined trajectory $\hat \tau$ is a sequence of features $(\hat f_{t+1},...,\hat f_{t+\tau})$, where $t$ is the current time, $\tau$ the length of rollout, and $\hat f_{t+i}$ the output of the environment model(the predicted observation and reward). In order to guarantee success in imperfections, the architecture does not assume the learned model to be perfect. The output will not only depend on the predicted reward.<br />
<br />
===Trajectories Encoder===<br />
From the intuition to keep the sequence information in the trajectories, the architecture uses a rollout encoder $\varepsilon$ that processes the imagined rollout as a whole and learns to interpret it(Figure 2 middle). Each trajectory is encoded as a rollout embedding $e_i=\varepsilon(\hat \tau_i)$. Then, the aggregator $A$ combines the rollout embedding s into a single imagination code $c_{ia}=A(e_1,...,e_n)$ by simply concatenating all the summaries.<br />
In the experiments, the encoder is an LSTM that takes the predicted output from environment model as the input. One observation is that the order of the sequence $\hat f_{t+1}$ to $\hat f_{t+\tau}$ makes relatively little impact on the performance. The encodes mimics the Bellman type backup operations in DQN.<br />
An alternative attempt would be to combine guided policy search with the linear quadratic regulator[7], which is coincidentally a joint model-free and model-based trajectory update mechanism for reinforcement learning.<br />
<br />
===Model-Free Path===<br />
The model-free path contains a network that only takes the current observation as input that generates the potential optimal action. This network can be same as the one in imagination core.<br />
<br />
<br />
In conclusion, the I2A learns to combine information for two paths, and without the model-based path, I2A simply reduce to a standard model-free network(such as A3C, more explanations [https://medium.com/emergent-future/simple-reinforcement-learning-with-tensorflow-part-8-asynchronous-actor-critic-agents-a3c-c88f72a5e9f2 here]). The imperfect approximation results in a rollout policy with higher entropy, potentially striking a balance between exploration and exploitation.<br />
<br />
=Experiment=<br />
These following experiments were tested in Sokoban and MiniPacman games. All results are averages taken from top three agents. These agents were trained over 32 to 64 workers, and the network was optimized by RMSprop.<br />
As the pre-training strategy, the training data of I2A was pre-generated from trajectories of a partially trained standard model-free agent, the data is also taken into account for the budget. The total number of frames that were needed in pre-training is counted in the later process. Meanwhile, the authors show that the environment model can be reused to solve multiple tasks in the same environment.<br />
<br />
In the game Sokoban, the environment is a 10 x 10 grid world. All agents were trained directly on raw pixels(image size 80 x 80 with 3 channels). To make sure the network is not just simply "memorize" all states, the game procedurally generates a new level each episode. Out of 40 million levels generated, less than 0.7% were repeated. Therefore, a good agent should solve the unseen level as well.<br />
<br />
The reward settings for reinforcement learning algorithms are as follows:<br />
* Every time step, a penalty of -0.1 is applied to the agent.(encourage agents to finish levels faster)<br />
* Whenever the agent pushes a box on target, it receives a reward of +1.(encourage agents to push boxes onto targets)<br />
* Whenever the agent pushes a box off target, it receives a penalty of -1.(avoid artificial reward loop that would be induced by repeatedly pushing a box off and on target)<br />
* Finishing the level gives the agent a reward of +10 and the level terminates.(strongly reward solving a level)<br />
<br />
To show the advantage of I2A, the authors set a model-free standard architecture as a baseline. The architecture is a multi-layer convolutional neural network (CNN), taking the current observation $O_t$ as input, followed by a fully connected (FC) hidden layer. This FC layer feeds into two heads: into an FC layer with one output per action computing the policy logits $\log \pi(a_t|O_t, \theta)$; and into another FC layer with a single output that computes the value function $V(O_t; \theta_v)$.<br />
* for MiniPacman: the CNN has two layers, both with 3x3 kernels, 16 output channels and strides 1 and 2; the following FC layer has 256 units<br />
* for Sokoban: the CNN has three layers with kernel sizes 8x8, 4x4, 3x3, strides of 4, 2, 1 and number of output channels 32, 64, 64; the following FC has 512 units<br />
<br />
===Sokoban===<br />
<br />
The environment model for Sokoban is shown in figure 4<br />
[[File:sokoban_em.png|400px|center|thumb|Figure 4: The Sokoban environment model]]<br />
<br />
Besides, to demonstrate the influence of larger architecture in I2A, the authors set a copy-model agent that uses the same architecture of I2A but the environment model is replaced by identical map. This agent is regarded as an I2A agent without imagination.<br />
<br />
[[File:sokoban_result.png|800px|center|thumb|Figure 5: Sokoban learning curves. Left: training curves of I2A and baselines. Right: I2A training curves for various values of imagination depth]]<br />
The results are shown in Figure 4(left). I2A agents can solve much more levels compared to common DQN. Also, it far outperforms the copy-model version, suggesting that the environment model is crucial. The authors also trained an I2A where the environment model was predicting no rewards, only observations. This also performed worse. However, after much longer training (3e9 steps), these agents did recover the performance of the original I2A, which was never the case for the baseline agent even with that many steps. Hence, reward prediction is very helpful but not absolutely necessary in this task, and imagined observations alone are informative enough to obtain high performance on Sokoban. Note this is in contrast to many classical planning and model-based reinforcement learning methods, which often rely on reward prediction.<br />
<br />
====Length of Rollout====<br />
A further experiment was investigating how the length of individual rollouts affects performance. The authors performed a parameter searching. Figure 5(right) shows the influence of the rollout length. The strategy using 3 rollout steps improves the speed of learning and improves the performance significantly than 1 step, and 5 is the optimal number. This implies rollout can be very helpful and informative. This rollout enables the agent to learn moves it cannot recover from.<br />
<br />
[[File:sokoban_noisy.png|800px|center|thumb|Figure 6: Experiments with a noisy environment model Left: each row shows an example 5-step rollout after conditioning on an environment observation. Errors accumulate and lead to various artifacts, including missing or duplicate sprites. Right: comparison of Monte-Carlo (MC) search and I2A when using either the accurate or the noisy model for rollouts.]]<br />
<br />
====Imperfections====<br />
To demonstrate I2A can handle less reliable predictions, the authors set experiment where the I2A used a poor environment model(smaller number of parameters), where the error may accumulate across the rollout(Figure 6 left). As shown in figure 6(right), even with relatively poor environment model, the performance of I2A is stable, unlike traditional Monte-Carlo search, which explicitly estimates the value of each action from rollouts, rather than learning an arbitrary encoding of the rollouts. An interesting result is that a rollout length 5 no longer outperforms a length of 3, which matches our common sense. The authors suggest that it is learning a rollout encoder that enables I2As to deal with imperfect model predictions.<br />
<br />
====Perfections====<br />
As I2A shows the robustness towards environment models, the authors tested an I2A agent with a nearly perfect environment model, and the results are in Table 1 and Table 2. Traditional Mento-Carlo Tree Search is tested as the baseline. From the table, although it is able to solve many levels, the search steps are very huge. On the contrary, I2A with the nearly perfect model can achieve the same fraction with much fewer steps.<br />
<br />
====Generalization====<br />
Lastly, the authors probe the generalization capabilities of I2As, beyond handling random level layouts in Sokoban. The agents were trained on levels with 4 boxes. Table 2 shows the performance of I2A when such an agent was tested on levels with different numbers of boxes, and that of the standard model-free agent for comparison. It turns out that I2As generalizes well; at 7 boxes, the I2A agent is still able to solve more than half of the levels, nearly as many as the standard agent on 4 boxes.<br />
[[File:i2a_table.png|800px|center|thumb]]<br />
<br />
===MiniPacman===<br />
MiniPacman is a game modified from the classical game PacMan. In the game(Figure 8, left), the player explores a maze that contains food while being chased by ghosts. The maze also contains power pills; when eaten, for a fixed number of steps, the player moves faster, and the ghosts run away and can be eaten. These dynamics are common to all tasks. Each task is defined by a vector $w \in R^5$, associating a reward to each of the following five events: moving, eating food, eating a power pill, eating a ghost, and being eaten by a ghost.<br />
The goal of this part is the attempt that tries to apply the same I2A model to different tasks. The five tasks are described as follows:<br />
* Regular: level is cleared when all the food is eaten;<br />
* Avoid: level is cleared after 128 steps;<br />
* Hunt: level is cleared when all ghosts are eaten or after 80 steps.<br />
* Ambush: level is cleared when all ghosts are eaten or after 80 steps.<br />
* Rush: level is cleared when all power pills are eaten.<br />
<br />
[[File:minipacman_reward.png|800px|center|thumb|Table 3: the reward settings in different tasks]]<br />
<br />
Different from the task in Sokoban, in order to capture long-range dependencies across pixels, the authors also made use of a layer that is called pool-and-inject, which applies global max-pooling over each feature map and broadcasts the resulting values as feature maps of the same size and concatenates the result to the input. Pool-and-inject layers are therefore size-preserving layers which communicate the max-value of each layer globally to the next convolutional layer. The environment model for MiniPacman is shown in Figure 7.<br />
<br />
[[File:minipacman_model.png|800px|center|thumb|Figure 7: The MiniPacman environment model]]<br />
<br />
To illustrate the benefits of model-based methods in this multi-task setting, the authors trained a single environment model to predict both observations (frames) and events, where the environment model is effectively shared across all tasks. Results in Figure 7(right) illustrates the benefit of the I2A architecture, outperforming the standard agent in all tasks. Note that for tasks 4 & 5, the rewards are particularly sparse, and the anticipation of ghost dynamics is especially important. The I2A agent can leverage its environment and reward model to explore the environment much more effectively.<br />
<br />
[[File:minipacman.png|800px|center|thumb|Figure 8: Minipacman environment Left: Two frames from a minipacman game: the player is green, dangerous ghosts red, food dark blue, empty corridors black, power pills in cyan. After eating a power pill (right frame), the player can eat the 4 weak ghosts (yellow). Right: Performance after 300 million environment steps for different agents and all tasks. Note I2A clearly outperforms the other two agents on all tasks with sparse rewards.]]<br />
<br />
=Conclusion=<br />
In this paper, the authors applied recent success in CNN and reinforcement learning and raised a novel approach, which is a combination of model-free and model-based methods, called Imagination-augmented RL. Unlike classical model-based RL and planning methods, I2A is able to successfully use imperfect models to support model-free decisions. This approach outperforms model-free baselines in Sokoban and MiniPacman games. As experiments suggest, this method is able to successfully use imperfect models to interpret future states and rewards. The imagination core part is essential in irreversible domains, where actions can have catastrophic outcomes. Compared to traditional Monte-Carlo search methods, the search space in I2A only grows linearly with the extension of the length of rollouts whereas I2As require far fewer function calls. This work may significantly broaden the applicability of model-based RL concepts and ideas.<br />
<br />
=Insight=<br />
This is a paper with very interesting ideas. However, it seems that the work is really hard to reproduce for an individual researcher. Since the architecture works as a whole, it is very difficult to debug each single part. Meanwhile, the training process is kind of long with up to 1e9 steps, which is also a huge requirement for computing resources.<br />
<br />
In terms of the architecture itself, the design the CNN for the tasks seems to be very empirical. The authors did not include the reasons or rules for this part. Yet why authors applied residual connection in this shadow network is unknown. According to the paper, even the CNN network is quite simple, some details in LSTM encoder are omitted. Therefore, the backpropagation process is not so clear across the whole model.<br />
<br />
Back to the settings of environment model, the authors used pre-trained model instead of the jointly training way. Would it be hard to train both models simultaneously?<br />
<br />
Lastly, the authors raised a new layer as Pool-and-inject layer, the motivation and plausibility are not so clear. It would be better if the authors can compare it with common pooling layer.<br />
<br />
In spite of some missing details, this is a solid work with a novel idea and many tricks. In addition, the settings of the experiment are quite inspiring where we can learn from.<br />
<br />
=Reference=<br />
# Buesing, L., Badia, A.P., Battaglia, P.W., Guez, A., Heess, N., Li, Y., Pascanu, R., Racanière, S., Reichert, D.P., Rezende, D.J., Silver, D., Vinyals, O., Weber, T., & Wierstra, D. (2017). Imagination-Augmented Agents for Deep Reinforcement Learning. CoRR, abs/1707.06203.<br />
# YuXuan Liu, Abhishek Gupta, Pieter Abbeel, and Sergey Levine. Imitation from observation: Learning to imitate behaviors from raw video via context translation. arXiv preprint arXiv:1707.03374, 2017.<br />
# Marc Deisenroth and Carl E Rasmussen. Pilco: A model-based and data-efficient approach to policy search. In Proceedings of the 28th International Conference on machine learning (ICML-11), pages 465–472, 2011.<br />
# Jessica B. Hamrick, Andy J. Ballard, Razvan Pascanu, Oriol Vinyals, Nicolas Heess, and Peter W. Battaglia. Metacontrol for adaptive imagination-based optimization. In Proceedings of the 5th International Conference on Learning Representations (ICLR 2017), 2017.<br />
# Razvan Pascanu, Yujia Li, Oriol Vinyals, Nicolas Heess, David Reichert, Theophane Weber, Sebastien Racaniere, Lars Buesing, Daan Wierstra, and Peter Battaglia. Learning model-based planning from scratch. arXiv preprint, 2017.<br />
# Mnih, V., Kavukcuoglu, K., Silver, D., Rusu, A. A., Veness, J., Bellemare, M. G., ... & Petersen, S. (2015). Human-level control through deep reinforcement learning. Nature, 518(7540), 529-533.<br />
#Yevgen Chebotar, Karol Hausman, Marvin Zhang, Gaurav Sukhatme, Stefan Schaal, Sergey Levine. "Combining Model-Based and Model-Free Updates for Trajectory-Centric Reinforcement Learning". arxiv pre-print; arXiv:1703.03078 [cs.RO]</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Learning_the_Number_of_Neurons_in_Deep_Networks&diff=29770Learning the Number of Neurons in Deep Networks2017-11-09T01:20:26Z<p>Nkadiyar: /* Model Training and Model Selection */</p>
<hr />
<div>='''Introduction'''=<br />
<br />
Due to the availability of large-scale datasets and powerful computation, '''Deep Learning''' has made huge breakthroughs in many areas, like Language Models and Computer Vision. In deep neural networks, we need to determine the number of layers and the number of neurons in each layer, i.e, we need to determine the number of parameters, or complexity of the model. Typically, this is determined by trial and error manually. Currently, this is mostly achieved by manually tuning these hyper-parameters using validation data or building very deep networks. However, building a very deep model is still challenging, especially for very large datasets, which leads to high cost on memory and reduction in speed.<br />
<br />
In this paper, we used an approach to automatically select the number of neurons in each layer when we learn the network. Our approach introduces a '''group sparsity regularizer''' on the parameters of the network, and each group acts on the parameters of one neuron, rather than training an initial network as as pre-processing step(training shallow or thin networks to mimic the behaviour of deep ones [Hinton et al., 2014, Romero et al., 2015]) and reducing neurons later as a post-processing step. We set those useless parameters to zero, which cancels out the effects of a particular neuron. Therefore, our approach does not need to learn a redundant network successfully and then reduce its parameters, instead, it learns the number of relevant neurons in each layer and the parameters of those neurons simultaneously.<br />
<br />
In the experiments on several image recognition datasets, we showed the effectiveness of our approach, which reduces the number of parameters by up to 80% compared to the complete model, and has no recognition accuracy loss at the same time. Actually, our approach even yields more effective and faster networks, and occupies less memory.<br />
<br />
='''Related Work'''=<br />
<br />
The recent researches tend to build very deep networks. Building very deep networks means we need to learn more parameters, which leads to a significant cost on the memory of the equipment as well as the speed. Even though automatic model selection has developed in the past years by constructive and destructive approaches, there are some drawbacks. For '''constructive method''', it starts a super shallow architecture, and then adds additional parameters [Bello, 1992]. A similar work that adds new layers to the initial shallow networks was successfully employed [Simonyan and Zisserman, 2014] at the process of learning. However, we know shallow networks have fewer parameters, so that it can not handle the non-linearities as effectively as the deep networks [Montufar et al., 2014], so shallow networks may easily get stuck by the bad optima. Therefore, the drawback of this method is that these networks may produce poor initializations for the later processes. The authors make this claim without ever providing any evidence for it. For '''destructive method''', it starts by a deep network to reduce a significant number of redundant parameters [Denil et al., 2013, Cheng et al., 2015] while keeping its behaviour unchanged. Even though this technique has shown removing the redundant parameters [LeCun et al., 1990, Hassibi et al., 1993] or the neurons [Mozer and Smolensky, 1988, Ji et al., 1990, Reed, 1993] has little influence on the output, it requires the analysis of each parameter and neuron by network Hessian, which is very computationally expensive for large architectures. The main motivation of these works was to build a more compact network.<br />
<br />
Particularly, building a compact network is a research focus for '''Convolutional Neural Networks'''(CNNs). Some works has proposed to decompose the filters of a pre-trained network into low-rank filters, which reduces the number of parameters [Jaderberg et al., 2014b, Denton et al., 2014, Gong et al., 2014]. The issue of this proposal is that we need to successfully train an initial deep network, since it acts as as post-processing step. [Weigend et al., 1991] and [Collins and Kohl, 2014] used direct training to develop regularizers that eliminate some of the parameters of the network. The problem is that the number of layers and neurons each layer is determined manually. A very similar work using the group lasso method for CNN was previously done in [Liu et al., 2015]. The big-picture idea appears to be very similar but they differ in details of methodology.<br />
<br />
='''Model Training and Model Selection'''=<br />
<br />
In general, a deep network has L layers containing linear operations on their inputs, intertwined with activation functions. The activation function we generally use is '''Rectified Linear Units(RELU) or sigmoids'''. Suppose each layer l has $N_{l}$ neurons, and each of them has parameters $\Theta=(\theta_{l})_{1\leqslant{l}\leqslant{L}}$, where $\theta_{l}=({\theta^n _{l}})_{1\leqslant{n}\leqslant{N_{l}}}$ and where $\theta^n _{l}=[w_{l}^{n},b_{l}^{n}]$ is consisted of linear component $w_{l}^{n}$ and linear bias $b_{l}^{n}$. Given an input $x$, under the linear, non-linear and pooling operations, we obtain the output $\hat{y}=f(x,\theta)$, where $f(*)$ encodes the succession of linear, non-linear and pooling operations.<br />
<br />
At the step of training, we have N input-output pairs ${(x_{i},y_{i})}_{1\leqslant{i}\leqslant{N}}$, and the loss function is given by $\ell(y_{i},f(x_{i},\Theta))$, which compares the predicted output with the ground-truth output. Generally, we choose logistic loss for classification and the square loss for regression. Therefore, learning the parameters of the network is equivalent to solving the optimization of the following:<br />
$$\displaystyle \min_{\Theta}\frac{1}{N}\sum_{i=1}^{N}\ell(y_{i},f(x_{i},\Theta))+\gamma(\Theta),$$ where $\gamma(\Theta)$ represents a regularizer on the network parameters. Our choice for the regularizer can be $\ell_{2}$-norm(i.e, weight decay) or $\ell_{1}$-norm. $\ell_{2}$-norm usually favours small parameter values, and $\ell_{1}$-norm can only delete those irrelevant parameters, but not the neurons. The goal in this paper is to automatically determine the number of neurons of each layer, but neither of the above techniques achieve this goal. Here, we make use of the '''group sparsity''' [Yuan and Lin., 2007] (starting from an overcomplete network and canceling the influence of some neurons). The regularizer, therefore, can be written as $$\gamma(\Theta)=\sum_{l=1}^{L}\beta_{l}\sqrt{P_{l}}\sum_{n=1}^{N_{l}}||\theta_{l}^{n}||_{2},$$ where $P_{l}$ means the size of the vector that includes the parameters of each neuron in layer $l$, and $\beta_{l}$ balances the influence of the penalty. In practice, we found the most effective way to select $\beta$ is a relatively small one for the first few layers, and a larger weight for the remaining layers. The reason we choose a small weight is that it can prevent deleting too much neurons in the first few layers, so that we have enough information for learning the remaining parameters. The original premise of this paper seemed to suggest a new method that was different from both the constructive and destructive methods described above. However, this approach of starting with an overcomplete network and training with group sparsity appears to be no different from destructive methods. The main contribution here is then the regularization function to act on entire neurons, which is in fairness an interesting approach.<br />
<br />
The group sparsity helps us effectively remove some of the neurons, and also standard regularizers on the individual parameters are effective for the generalization purpose [Bartlett, 19996, Krogh and Hertz, 1992, Theodoridis, 2015, Collins and Kohli, 2014]. By this idea, we introduce '''sparse group Lasso''', which considers a more generalised penalty that merges L1 norm in Lasso with the group lasso (i.e. "two-norm"). This leads to the production of a penalty which specifies solutions that are sparse enough both at an individual and group feature levels [1]. It specifies that the regularizer can be written as $$\gamma(\Theta)=\sum_{l=1}^{L}((1-\alpha)\beta_{l}\sqrt{P_{l}}\sum_{n=1}^{N_{l}}||\theta_{l}^{n}||_{2}+\alpha\beta_{l}||\theta_{l}||_{1},$$ where $\alpha\in[0,1]$. We find that if $\alpha=0$, then we have the group sparsity regularizer. In practice, we use both $\alpha=0$ and $\alpha=0.5$ in the experiments.<br />
<br />
This reminds me of the relationships among Lasso regression, Ridge regression and Elastic Net regression. In lasso regression, the penalized residual sum of squares is composed of the regular residual sum of squared plus a L1 regularizer. In ridge regression, its penalized residual sum of squares is composed of the regular residual sum of squared plus a L2 regularizer. Finally, an elastic net regression is a combination of lasso regularizer and ridge regularizer, where its objective function is to optimize parameters by including both L1 and L2 norms. <br />
<br />
To find the optimization, in this paper we use proximal gradient descent [Parikh and Boyed, 2014]. This approach iteratively takes a gradient step of size t with respect to the loss. The following is the algorithm for it: <br />
<br />
We define proximal operator of f as $$prox_{f}(v)=\displaystyle \min_{x}(\frac{1}{2t}||x-v||_{2}^{2}+f(x))$$ <br />
<br />
<br />
Suppose we want to minimize $f(x)+g(x)$, and the proximal gradient method is given by $$x^{(k+1)}=prox_{t^{k}g}(x^{k}-t^{k}\nabla{f}(x^{k})), k=1,2,3...$$ <br />
<br />
Therefore, we can update our parameter by the above method as $$\tilde{\theta}_{l}^{n}=\displaystyle \min_{\theta_{l}^{n}}\frac{1}{2t}||\theta_{l}^{n}-\hat{\theta}_{l}^{n}||_{2}^{2}+\gamma(\Theta),$$<br />
where $\hat{\theta}_{l}^{n}$ is the solution obtained from the general loss gradient. By the derivative of [Simon et al., 2013], we have a closed-form solution for this problem: <br />
$$\tilde{\theta}_{l}^{n}=(1-\frac{t(1-\alpha)\beta_{l}\sqrt{P_{l}}}{||S(\hat{\theta}_{l}^{n},t\alpha\beta_{l})||_{2})})_{+}S(\hat{\theta}_{l}^{n},t\alpha\beta_{l}),$$<br />
where + refers to taking the maximum between the argument and 0, and $S(*)$ is $$S(a,b)=sign(a)(|a|-b)_{+}$$<br />
In practice, we use stochastic gradient descent and work with mini-batches, and then update the variables of all the groups according to the closed-form of $\tilde{\theta}_{l}^{n}$. When the learning steps terminate, we remove the neurons whose parameters have gone to zero. Additionally, when examining fully-connected layers, the neurons acting on output of zeroed-out neurons of the previous layer also become useless, and are removed accordingly.<br />
<br />
='''Experiment'''=<br />
<br />
==='''Set Up'''===<br />
<br />
They use two large-scale image classification datasets, '''ImageNet''' [Russakovsky et al., 2015] and '''Places2-401''' [Zhou et al., 2015]. They also conducted additional experiments on the '''ICDAR''' character recognition dataset of [Jaderberg et al., 2014a]. <br />
<br />
For ImageNet, they used the subset which contains 1000 categories, with 1.2 million training images and 50000 validation images. For Places2-401, it has more than 10 million images with 401 unique scene categories. 5000 to 30000 images are comprised into per category. Both architectures of these two datasets are based on the VGG-B network(BNet) [Simonyan and Zisserman, 2014] and on DecomposeMe8($Dec_{8}$) [ALvarez and Petersson, 2016]. BNet has 10 convolutional layers followed by 3 fully-connected layers. In the experiment, they remove the first 2 fully-connected layers, which we call $BNet^{C}$. $Dec_{8}$ contains 16 convolutional layers with 1D kernels, which can model 8 2D convolutional layers. Both models were trained for a total of 55 epochs with 12000 batches per epoch and a batch size of 48 and 180 for BNet and $Dec_{8}$, respectively. The learning rate was initialized by 0.01 and then multiplied by 0.1. They set $\beta_{l}$=0.102 for the first three layers and $\beta_{l}$=0.255 for the remaining ones.<br />
<br />
For ICDAR dataset, it consists of 185639 training and 5198 test data split into 36 categories. The architecture here starts 6 1D convolutional layers with max-pooling, rather than 3 convolutional layers with a maxout layer [Goodfellow et al., 2013] after each convolution, followed by one fully-connected layer. They call their architecture as Dec3. The model was trained for a total of 45 epochs with a batch size of 256 and 1000 iterations per epoch. The learning rate was initialized by 0.1 and multiplied by 0.1 in the second, seventh and fifteenth epochs. They set $\beta_{l}$=5.1 for the first layer and $\beta_{l}$=10.2 for the remaining ones.<br />
<br />
==='''Results'''===<br />
<br />
[[File:imageNet.png]]<br />
<br />
The above table show the accuracy comparisons between the original architectures and ours. For $Dec_{8}$ on the ImageNet dataset, we evaluated two additional models: $Dec_{8}-640$ with 640 neurons per layer and $Dec_{8}-768$ with 768 neurons per layer. $Dec_{8}-640_{SGL}$ means the sparse group Lasso regularizer with $\alpha=0.5$ and $Dec_{8}-640_{GS}$ represents the group sparsity regularizer. Note that all our architectures yield an improvement over the original network except $Dec_{8}-768$. For instance, Ours-$Bnet_{GS}^{C}$ increases the performance of 1.6% compared to $BNet^{C}$. <br />
<br />
[[File:44.png]]<br />
<br />
[[File:2.png]]<br />
<br />
The above figures reports the reduced percentage of neurons/parameters with our approach for $BNet^{C}$ and $Dec_{8}$. For example, in the first figure, our approach reduces the number of neurons by over 12% and the number of parameters by around 14%, while improving the generalization ability of 1.6%(as indicated by accuracy gap). The left image in the first figure also shows that reduction in number of neurons is spread all the layers with the largest difference in the L10. For $Dec_{8}$, in the second figure, we find when we increase the number of neurons in each layer, the benefits of our approach become more significant. For instance, $Dec_{8}-640$ with group sparsity regularizer reduces the number of neurons by 10%, and of parameters by 12.48%. The left image in the second figure also shows that reduction in number of neurons is spread all the layers. <br />
<br />
[[File:ICDA.png]]<br />
<br />
Finally, the above figure indicates the experiment results for ICDAR dataset. Here, we used the $Dec_{3}$ architecture, where the last two layers initially contain 512 neurons. The accuracy rate for $MaxPllo_{2Dneurons}$ is 83.8%, and accuracy rate for $Dec_{3}$ is 89.3%, which means 1D filters perform better than a network with 2D kernels. Our model on this dataset reduces 38.64% of neurons and totally up to 80% of the number of parameters with group sparsity regularizer.<br />
<br />
All the above results evidence that our algorithm effectively removes the number of parameters and increases the model accuracy. Our algorithm of automatic model selection effectively performs on the classification task.<br />
<br />
='''Analysis on Testing'''=<br />
<br />
Our algorithm does not remove neurons during the training time, however, we remove those neurons after training, which yields a smaller network at test time. This improvement not only reduces the number of parameters of the network, but also decreases the computational memory cost and increases the speed. <br />
<br />
[[File:table2.png]]<br />
<br />
The above table reports the runtime, memory, as well as the percentage of reduced parameters after removing the zeroed-out neurons. The BNet and $Dec_{8}$ were tested on the dataset of ImageNet, while $Dec_{3-GS}$ was tested on the dataset of ICDAR. From the table, we find that all the models for the ImageNet and ICDAR have speeded up the runtime, for example, $Dec_{8}-768_{GS}$ on ImageNet data speeds up the runtime nearly 16% at the batch size of 8, and $Dec_{3}$ on ICDAR data speeds up nearly 50% at natch size of 16. For the percentage of parameters reduced, we find BNet, $Dec_{8}-640_{GS}$ and $Dec_{8}-768_{GS}$ have reduced 12.06%, 26.51%, and 46.73% respectively. More significantly, for $Dec_{3-GS}$, it reduces 82.35% of the parameters. All of these changes show the benefits at the testing time.<br />
<br />
='''Conclusion'''=<br />
<br />
In this paper, we have introduced an approach that relies on group sparsity regularizer. This approach automatically determines the number of neurons in each layer of a deep network. From the experiments, we found our approach not only reduces the number of parameters in our model, but also saves the computation memory and increases the speed at test time. However, the limitation of our approach is that the number of layers in the network remains fixed.<br />
<br />
='''Critique'''=<br />
The authors of the paper state that "...we assume that the parameters of each neuron in layer $l$ are grouped in a vector of size $P_{l}$ and where $\lambda_{l}$ sets the influence of the penalty. Note that, in the general case, this weight can be different for each layer $l$. In practice, however, we found most effective to have<br />
two different weights: a relatively small one for the first few layers, and a larger weight for the<br />
remaining ones. This effectively prevents killing too many neurons in the first few layers, and thus<br />
retains enough information for the remaining ones." However, the authors fail to present any guidance as to what gets counted as "the first few layers" and what the relative sizes for the two weights should be even after we have chosen the "first few layers". Indeed, such choice seems to be an unaccounted component of tuning the model but this receives scant attention in the current paper. Several numerical comparisons should be carried out to allow further discussion on this question.<br />
<br />
The experiments could have included better baseline models to compare against. For example, how do we know the original model was not overly complex to begin with? It might have been a good idea for the authors to compare their group sparse lasso method against the naive method of (blindly) reducing the number of neurons in each layer by 10-20% just for a very preliminary check. Also, the author could have applied the theory of ridge and Lasso regression to analyze the effect of the regularization mathematically.<br />
<br />
A rather reliable method of experimentation to compare the performance and accuracy has been left out. The authors have not stated any comparisons of this method with the Dropout method [Srivastava,2014], which is similar in terms of the physical effects on the network.<br />
<br />
='''References'''=<br />
<br />
P. L. Bartlett. For valid generalization the size of the weights is more important than the size of the network. In NIPS, 1996.<br />
<br />
M. G. Bello. Enhanced training algorithms, and integrated training/architecture selection for multilayer perceptron networks. IEEE Transactions on Neural Networks, 3(6):864–875, Nov 1992.<br />
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Yu Cheng, Felix X. Yu, Rogério Schmidt Feris, Sanjiv Kumar, Alok N. Choudhary, and Shih-Fu Chang. An exploration of parameter redundancy in deep networks with circulant projections. In ICCV, 2015.<br />
<br />
I. J. Goodfellow, D. Warde-farley, M. Mirza, A. Courville, and Y. Bengio. Maxout networks. In ICML, 2013.<br />
<br />
G. E. Hinton, O. Vinyals, and J. Dean. Distilling the knowledge in a neural network. In arXiv, 2014.<br />
<br />
M. Jaderberg, A. Vedaldi, and A. Zisserman. Deep features for text spotting. In ECCV, 2014a.<br />
<br />
M. Jaderberg, A. Vedaldi, and A. Zisserman. Speeding up convolutional neural networks with low rank expansions. In BMVC, 2014b.<br />
<br />
N. Simon, J. Friedman, T. Hastie, and R. Tibshirani. A sparse-group lasso. Journal of Computational and Graphical Statistics, 2013.<br />
<br />
H. Zhou, J. M. Alvarez, and F. Porikli. Less is more: Towards compact CNNs. In ECCV, 2016.<br />
<br />
Group LASSO - https://pdfs.semanticscholar.org/f677/a011b2a912e3c5c604f6872b9716cc0b8aa0.pdf<br />
<br />
Liu, Baoyuan, et al. "Sparse convolutional neural networks." Proceedings of the IEEE Conference on Computer Vision and Pattern Recognition. 2015.<br />
<br />
Nitish Srivastava, Geoffrey Hinton, Alex Krizhevsky, Ilya Sutskever, and Ruslan Salakhutdinov. 2014. Dropout: a simple way to prevent neural networks from overfitting. J. Mach. Learn. Res. 15, 1 (January 2014), 1929-1958.<br />
<br />
<br />
Derivation & Motivation of the Soft Thresholding Operator (Proximal Operator):<br />
# http://www.onmyphd.com/?p=proximal.operator<br />
# https://math.stackexchange.com/questions/471339/derivation-of-soft-thresholding-operator</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Learning_the_Number_of_Neurons_in_Deep_Networks&diff=29768Learning the Number of Neurons in Deep Networks2017-11-09T01:11:25Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>='''Introduction'''=<br />
<br />
Due to the availability of large-scale datasets and powerful computation, '''Deep Learning''' has made huge breakthroughs in many areas, like Language Models and Computer Vision. In deep neural networks, we need to determine the number of layers and the number of neurons in each layer, i.e, we need to determine the number of parameters, or complexity of the model. Typically, this is determined by trial and error manually. Currently, this is mostly achieved by manually tuning these hyper-parameters using validation data or building very deep networks. However, building a very deep model is still challenging, especially for very large datasets, which leads to high cost on memory and reduction in speed.<br />
<br />
In this paper, we used an approach to automatically select the number of neurons in each layer when we learn the network. Our approach introduces a '''group sparsity regularizer''' on the parameters of the network, and each group acts on the parameters of one neuron, rather than training an initial network as as pre-processing step(training shallow or thin networks to mimic the behaviour of deep ones [Hinton et al., 2014, Romero et al., 2015]) and reducing neurons later as a post-processing step. We set those useless parameters to zero, which cancels out the effects of a particular neuron. Therefore, our approach does not need to learn a redundant network successfully and then reduce its parameters, instead, it learns the number of relevant neurons in each layer and the parameters of those neurons simultaneously.<br />
<br />
In the experiments on several image recognition datasets, we showed the effectiveness of our approach, which reduces the number of parameters by up to 80% compared to the complete model, and has no recognition accuracy loss at the same time. Actually, our approach even yields more effective and faster networks, and occupies less memory.<br />
<br />
='''Related Work'''=<br />
<br />
The recent researches tend to build very deep networks. Building very deep networks means we need to learn more parameters, which leads to a significant cost on the memory of the equipment as well as the speed. Even though automatic model selection has developed in the past years by constructive and destructive approaches, there are some drawbacks. For '''constructive method''', it starts a super shallow architecture, and then adds additional parameters [Bello, 1992]. A similar work that adds new layers to the initial shallow networks was successfully employed [Simonyan and Zisserman, 2014] at the process of learning. However, we know shallow networks have fewer parameters, so that it can not handle the non-linearities as effectively as the deep networks [Montufar et al., 2014], so shallow networks may easily get stuck by the bad optima. Therefore, the drawback of this method is that these networks may produce poor initializations for the later processes. The authors make this claim without ever providing any evidence for it. For '''destructive method''', it starts by a deep network to reduce a significant number of redundant parameters [Denil et al., 2013, Cheng et al., 2015] while keeping its behaviour unchanged. Even though this technique has shown removing the redundant parameters [LeCun et al., 1990, Hassibi et al., 1993] or the neurons [Mozer and Smolensky, 1988, Ji et al., 1990, Reed, 1993] has little influence on the output, it requires the analysis of each parameter and neuron by network Hessian, which is very computationally expensive for large architectures. The main motivation of these works was to build a more compact network.<br />
<br />
Particularly, building a compact network is a research focus for '''Convolutional Neural Networks'''(CNNs). Some works has proposed to decompose the filters of a pre-trained network into low-rank filters, which reduces the number of parameters [Jaderberg et al., 2014b, Denton et al., 2014, Gong et al., 2014]. The issue of this proposal is that we need to successfully train an initial deep network, since it acts as as post-processing step. [Weigend et al., 1991] and [Collins and Kohl, 2014] used direct training to develop regularizers that eliminate some of the parameters of the network. The problem is that the number of layers and neurons each layer is determined manually. A very similar work using the group lasso method for CNN was previously done in [Liu et al., 2015]. The big-picture idea appears to be very similar but they differ in details of methodology.<br />
<br />
='''Model Training and Model Selection'''=<br />
<br />
In general, a deep network has L layers containing linear operations on their inputs, intertwined with activation functions. The activation function we generally use is '''Rectified Linear Units(RELU) or sigmoids'''. Suppose each layer l has $N_{l}$ neurons, and each of them has parameters $\Theta=(\theta_{l})_{1\leqslant{l}\leqslant{L}}$, where $\theta_{l}=({\theta^n _{l}})_{1\leqslant{n}\leqslant{N_{l}}}$ and where $\theta^n _{l}=[w_{l}^{n},b_{l}^{n}]$ is consisted of linear component $w_{l}^{n}$ and linear bias $b_{l}^{n}$. Given an input $x$, under the linear, on-linear and pooling operations, we obtain the output $\hat{y}=f(x,\theta)$, where $f(*)$ encodes the succession of linear, non-linear and pooling operations.<br />
<br />
At the step of training, we have N input-output pairs ${(x_{i},y_{i})}_{1\leqslant{i}\leqslant{N}}$, and the loss function is given by $\ell(y_{i},f(x_{i},\Theta))$, which compares the predicted output with the ground-truth output. Generally, we choose logistic loss for classification and the square loss for regression. Therefore, learning the parameters of the network is equivalent to solving the optimization of the following:<br />
$$\displaystyle \min_{\Theta}\frac{1}{N}\sum_{i=1}^{N}\ell(y_{i},f(x_{i},\Theta))+\gamma(\Theta),$$ where $\gamma(\Theta)$ represents a regularizer on the network parameters. Our choice for the regularizer can be $\ell_{2}$-norm(i.e, weight decay) or $\ell_{1}$-norm. $\ell_{2}$-norm usually favours small parameter values, and $\ell_{1}$-norm can only delete those irrelevant parameters, but not the neurons. The goal in this paper is to automatically determine the number of neurons of each layer, but neither of the above techniques achieve this goal. Here, we make use of the '''group sparsity''' [Yuan and Lin., 2007] (starting from an overcomplete network and canceling the influence of some neurons). The regularizer, therefore, can be written as $$\gamma(\Theta)=\sum_{l=1}^{L}\beta_{l}\sqrt{P_{l}}\sum_{n=1}^{N_{l}}||\theta_{l}^{n}||_{2},$$ where $P_{l}$ means the size of the vector that includes the parameters of each neuron in layer $l$, and $\beta_{l}$ balances the influence of the penalty. In practice, we found the most effective way to select $\beta$ is a relatively small one for the first few layers, and a larger weight for the remaining layers. The reason we choose a small weight is that it can prevent deleting too much neurons in the first few layers, so that we have enough information for learning the remaining parameters. The original premise of this paper seemed to suggest a new method that was different from both the constructive and destructive methods described above. However, this approach of starting with an overcomplete network and training with group sparsity appears to be no different from destructive methods. The main contribution here is then the regularization function to act on entire neurons, which is in fairness an interesting approach.<br />
<br />
The group sparsity helps us effectively remove some of the neurons, and also standard regularizers on the individual parameters are effective for the generalization purpose [Bartlett, 19996, Krogh and Hertz, 1992, Theodoridis, 2015, Collins and Kohli, 2014]. By this idea, we introduce '''sparse group Lasso''', which considers a more generalised penalty that merges L1 norm in Lasso with the group lasso (i.e. "two-norm"). This leads to the production of a penalty which specifies solutions that are sparse enough both at an individual and group feature levels [1]. It specifies that the regularizer can be written as $$\gamma(\Theta)=\sum_{l=1}^{L}((1-\alpha)\beta_{l}\sqrt{P_{l}}\sum_{n=1}^{N_{l}}||\theta_{l}^{n}||_{2}+\alpha\beta_{l}||\theta_{l}||_{1},$$ where $\alpha\in[0,1]$. We find that if $\alpha=0$, then we have the group sparsity regularizer. In practice, we use both $\alpha=0$ and $\alpha=0.5$ in the experiments.<br />
<br />
This reminds me of the relationships among Lasso regression, Ridge regression and Elastic Net regression. In lasso regression, the penalized residual sum of squares is composed of the regular residual sum of squared plus a L1 regularizer. In ridge regression, its penalized residual sum of squares is composed of the regular residual sum of squared plus a L2 regularizer. Finally, an elastic net regression is a combination of lasso regularizer and ridge regularizer, where its objective function is to optimize parameters by including both L1 and L2 norms. <br />
<br />
To find the optimization, in this paper we use proximal gradient descent [Parikh and Boyed, 2014]. This approach iteratively takes a gradient step of size t with respect to the loss. The following is the algorithm for it: <br />
<br />
We define proximal operator of f as $$prox_{f}(v)=\displaystyle \min_{x}(\frac{1}{2t}||x-v||_{2}^{2}+f(x))$$ <br />
<br />
<br />
Suppose we want to minimize $f(x)+g(x)$, and the proximal gradient method is given by $$x^{(k+1)}=prox_{t^{k}g}(x^{k}-t^{k}\nabla{f}(x^{k})), k=1,2,3...$$ <br />
<br />
Therefore, we can update our parameter by the above method as $$\tilde{\theta}_{l}^{n}=\displaystyle \min_{\theta_{l}^{n}}\frac{1}{2t}||\theta_{l}^{n}-\hat{\theta}_{l}^{n}||_{2}^{2}+\gamma(\Theta),$$<br />
where $\hat{\theta}_{l}^{n}$ is the solution obtained from the general loss gradient. By the derivative of [Simon et al., 2013], we have a closed-form solution for this problem: <br />
$$\tilde{\theta}_{l}^{n}=(1-\frac{t(1-\alpha)\beta_{l}\sqrt{P_{l}}}{||S(\hat{\theta}_{l}^{n},t\alpha\beta_{l})||_{2})})_{+}S(\hat{\theta}_{l}^{n},t\alpha\beta_{l}),$$<br />
where + refers to taking the maximum between the argument and 0, and $S(*)$ is $$S(a,b)=sign(a)(|a|-b)_{+}$$<br />
In practice, we use stochastic gradient descent and work with mini-batches, and then update the variables of all the groups according to the closed-form of $\tilde{\theta}_{l}^{n}$. When the learning steps terminate, we remove the neurons whose parameters have gone to zero. Additionally, when examining fully-connected layers, the neurons acting on output of zeroed-out neurons of the previous layer also become useless, and are removed accordingly.<br />
<br />
='''Experiment'''=<br />
<br />
==='''Set Up'''===<br />
<br />
They use two large-scale image classification datasets, '''ImageNet''' [Russakovsky et al., 2015] and '''Places2-401''' [Zhou et al., 2015]. They also conducted additional experiments on the '''ICDAR''' character recognition dataset of [Jaderberg et al., 2014a]. <br />
<br />
For ImageNet, they used the subset which contains 1000 categories, with 1.2 million training images and 50000 validation images. For Places2-401, it has more than 10 million images with 401 unique scene categories. 5000 to 30000 images are comprised into per category. Both architectures of these two datasets are based on the VGG-B network(BNet) [Simonyan and Zisserman, 2014] and on DecomposeMe8($Dec_{8}$) [ALvarez and Petersson, 2016]. BNet has 10 convolutional layers followed by 3 fully-connected layers. In the experiment, they remove the first 2 fully-connected layers, which we call $BNet^{C}$. $Dec_{8}$ contains 16 convolutional layers with 1D kernels, which can model 8 2D convolutional layers. Both models were trained for a total of 55 epochs with 12000 batches per epoch and a batch size of 48 and 180 for BNet and $Dec_{8}$, respectively. The learning rate was initialized by 0.01 and then multiplied by 0.1. They set $\beta_{l}$=0.102 for the first three layers and $\beta_{l}$=0.255 for the remaining ones.<br />
<br />
For ICDAR dataset, it consists of 185639 training and 5198 test data split into 36 categories. The architecture here starts 6 1D convolutional layers with max-pooling, rather than 3 convolutional layers with a maxout layer [Goodfellow et al., 2013] after each convolution, followed by one fully-connected layer. They call their architecture as Dec3. The model was trained for a total of 45 epochs with a batch size of 256 and 1000 iterations per epoch. The learning rate was initialized by 0.1 and multiplied by 0.1 in the second, seventh and fifteenth epochs. They set $\beta_{l}$=5.1 for the first layer and $\beta_{l}$=10.2 for the remaining ones.<br />
<br />
==='''Results'''===<br />
<br />
[[File:imageNet.png]]<br />
<br />
The above table show the accuracy comparisons between the original architectures and ours. For $Dec_{8}$ on the ImageNet dataset, we evaluated two additional models: $Dec_{8}-640$ with 640 neurons per layer and $Dec_{8}-768$ with 768 neurons per layer. $Dec_{8}-640_{SGL}$ means the sparse group Lasso regularizer with $\alpha=0.5$ and $Dec_{8}-640_{GS}$ represents the group sparsity regularizer. Note that all our architectures yield an improvement over the original network except $Dec_{8}-768$. For instance, Ours-$Bnet_{GS}^{C}$ increases the performance of 1.6% compared to $BNet^{C}$. <br />
<br />
[[File:44.png]]<br />
<br />
[[File:2.png]]<br />
<br />
The above figures reports the reduced percentage of neurons/parameters with our approach for $BNet^{C}$ and $Dec_{8}$. For example, in the first figure, our approach reduces the number of neurons by over 12% and the number of parameters by around 14%, while improving the generalization ability of 1.6%(as indicated by accuracy gap). The left image in the first figure also shows that reduction in number of neurons is spread all the layers with the largest difference in the L10. For $Dec_{8}$, in the second figure, we find when we increase the number of neurons in each layer, the benefits of our approach become more significant. For instance, $Dec_{8}-640$ with group sparsity regularizer reduces the number of neurons by 10%, and of parameters by 12.48%. The left image in the second figure also shows that reduction in number of neurons is spread all the layers. <br />
<br />
[[File:ICDA.png]]<br />
<br />
Finally, the above figure indicates the experiment results for ICDAR dataset. Here, we used the $Dec_{3}$ architecture, where the last two layers initially contain 512 neurons. The accuracy rate for $MaxPllo_{2Dneurons}$ is 83.8%, and accuracy rate for $Dec_{3}$ is 89.3%, which means 1D filters perform better than a network with 2D kernels. Our model on this dataset reduces 38.64% of neurons and totally up to 80% of the number of parameters with group sparsity regularizer.<br />
<br />
All the above results evidence that our algorithm effectively removes the number of parameters and increases the model accuracy. Our algorithm of automatic model selection effectively performs on the classification task.<br />
<br />
='''Analysis on Testing'''=<br />
<br />
Our algorithm does not remove neurons during the training time, however, we remove those neurons after training, which yields a smaller network at test time. This improvement not only reduces the number of parameters of the network, but also decreases the computational memory cost and increases the speed. <br />
<br />
[[File:table2.png]]<br />
<br />
The above table reports the runtime, memory, as well as the percentage of reduced parameters after removing the zeroed-out neurons. The BNet and $Dec_{8}$ were tested on the dataset of ImageNet, while $Dec_{3-GS}$ was tested on the dataset of ICDAR. From the table, we find that all the models for the ImageNet and ICDAR have speeded up the runtime, for example, $Dec_{8}-768_{GS}$ on ImageNet data speeds up the runtime nearly 16% at the batch size of 8, and $Dec_{3}$ on ICDAR data speeds up nearly 50% at natch size of 16. For the percentage of parameters reduced, we find BNet, $Dec_{8}-640_{GS}$ and $Dec_{8}-768_{GS}$ have reduced 12.06%, 26.51%, and 46.73% respectively. More significantly, for $Dec_{3-GS}$, it reduces 82.35% of the parameters. All of these changes show the benefits at the testing time.<br />
<br />
='''Conclusion'''=<br />
<br />
In this paper, we have introduced an approach that relies on group sparsity regularizer. This approach automatically determines the number of neurons in each layer of a deep network. From the experiments, we found our approach not only reduces the number of parameters in our model, but also saves the computation memory and increases the speed at test time. However, the limitation of our approach is that the number of layers in the network remains fixed.<br />
<br />
='''Critique'''=<br />
The authors of the paper state that "...we assume that the parameters of each neuron in layer $l$ are grouped in a vector of size $P_{l}$ and where $\lambda_{l}$ sets the influence of the penalty. Note that, in the general case, this weight can be different for each layer $l$. In practice, however, we found most effective to have<br />
two different weights: a relatively small one for the first few layers, and a larger weight for the<br />
remaining ones. This effectively prevents killing too many neurons in the first few layers, and thus<br />
retains enough information for the remaining ones." However, the authors fail to present any guidance as to what gets counted as "the first few layers" and what the relative sizes for the two weights should be even after we have chosen the "first few layers". Indeed, such choice seems to be an unaccounted component of tuning the model but this receives scant attention in the current paper. Several numerical comparisons should be carried out to allow further discussion on this question.<br />
<br />
The experiments could have included better baseline models to compare against. For example, how do we know the original model was not overly complex to begin with? It might have been a good idea for the authors to compare their group sparse lasso method against the naive method of (blindly) reducing the number of neurons in each layer by 10-20% just for a very preliminary check. Also, the author could have applied the theory of ridge and Lasso regression to analyze the effect of the regularization mathematically.<br />
<br />
A rather reliable method of experimentation to compare the performance and accuracy has been left out. The authors have not stated any comparisons of this method with the Dropout method [Srivastava,2014], which is similar in terms of the physical effects on the network.<br />
<br />
='''References'''=<br />
<br />
P. L. Bartlett. For valid generalization the size of the weights is more important than the size of the network. In NIPS, 1996.<br />
<br />
M. G. Bello. Enhanced training algorithms, and integrated training/architecture selection for multilayer perceptron networks. IEEE Transactions on Neural Networks, 3(6):864–875, Nov 1992.<br />
<br />
Yu Cheng, Felix X. Yu, Rogério Schmidt Feris, Sanjiv Kumar, Alok N. Choudhary, and Shih-Fu Chang. An exploration of parameter redundancy in deep networks with circulant projections. In ICCV, 2015.<br />
<br />
I. J. Goodfellow, D. Warde-farley, M. Mirza, A. Courville, and Y. Bengio. Maxout networks. In ICML, 2013.<br />
<br />
G. E. Hinton, O. Vinyals, and J. Dean. Distilling the knowledge in a neural network. In arXiv, 2014.<br />
<br />
M. Jaderberg, A. Vedaldi, and A. Zisserman. Deep features for text spotting. In ECCV, 2014a.<br />
<br />
M. Jaderberg, A. Vedaldi, and A. Zisserman. Speeding up convolutional neural networks with low rank expansions. In BMVC, 2014b.<br />
<br />
N. Simon, J. Friedman, T. Hastie, and R. Tibshirani. A sparse-group lasso. Journal of Computational and Graphical Statistics, 2013.<br />
<br />
H. Zhou, J. M. Alvarez, and F. Porikli. Less is more: Towards compact CNNs. In ECCV, 2016.<br />
<br />
Group LASSO - https://pdfs.semanticscholar.org/f677/a011b2a912e3c5c604f6872b9716cc0b8aa0.pdf<br />
<br />
Liu, Baoyuan, et al. "Sparse convolutional neural networks." Proceedings of the IEEE Conference on Computer Vision and Pattern Recognition. 2015.<br />
<br />
Nitish Srivastava, Geoffrey Hinton, Alex Krizhevsky, Ilya Sutskever, and Ruslan Salakhutdinov. 2014. Dropout: a simple way to prevent neural networks from overfitting. J. Mach. Learn. Res. 15, 1 (January 2014), 1929-1958.<br />
<br />
<br />
Derivation & Motivation of the Soft Thresholding Operator (Proximal Operator):<br />
# http://www.onmyphd.com/?p=proximal.operator<br />
# https://math.stackexchange.com/questions/471339/derivation-of-soft-thresholding-operator</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Learning_the_Number_of_Neurons_in_Deep_Networks&diff=29767Learning the Number of Neurons in Deep Networks2017-11-09T01:08:39Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>='''Introduction'''=<br />
<br />
Due to the availability of large-scale datasets and powerful computation, '''Deep Learning''' has made huge breakthroughs in many areas, like Language Models and Computer Vision. In deep neural networks, we need to determine the number of layers and the number of neurons in each layer, i.e, we need to determine the number of parameters, or complexity of the model. Typically, this is determined by trial and error manually. Currently, this is mostly achieved by manually tuning these hyper-parameters using validation data or building very deep networks. However, building a very deep model is still challenging, especially for very large datasets, which leads to high cost on memory and reduction in speed.<br />
<br />
In this paper, we used an approach to automatically select the number of neurons in each layer when we learn the network. Our approach introduces a '''group sparsity regularizer''' on the parameters of the network, and each group acts on the parameters of one neuron, rather than trains an initial network as as pre-processing step(training shallow or thin networks to mimic the behaviour of deep ones [Hinton et al., 2014, Romero et al., 2015]). We set those useless parameters to zero, which cancels out the effects of a particular neuron. Therefore, our approach does not need to learn a redundant network successfully and then reduce its parameters, instead, it learns the number of relevant neurons in each layer and the parameters of those neurons simultaneously.<br />
<br />
In the experiments on several image recognition datasets, we showed the effectiveness of our approach, which reduces the number of parameters by up to 80% compared to the complete model, and has no recognition accuracy loss at the same time. Actually, our approach even yields more effective and faster networks, and occupies less memory.<br />
<br />
='''Related Work'''=<br />
<br />
The recent researches tend to build very deep networks. Building very deep networks means we need to learn more parameters, which leads to a significant cost on the memory of the equipment as well as the speed. Even though automatic model selection has developed in the past years by constructive and destructive approaches, there are some drawbacks. For '''constructive method''', it starts a super shallow architecture, and then adds additional parameters [Bello, 1992]. A similar work that adds new layers to the initial shallow networks was successfully employed [Simonyan and Zisserman, 2014] at the process of learning. However, we know shallow networks have fewer parameters, so that it can not handle the non-linearities as effectively as the deep networks [Montufar et al., 2014], so shallow networks may easily get stuck by the bad optima. Therefore, the drawback of this method is that these networks may produce poor initializations for the later processes. The authors make this claim without ever providing any evidence for it. For '''destructive method''', it starts by a deep network to reduce a significant number of redundant parameters [Denil et al., 2013, Cheng et al., 2015] while keeping its behaviour unchanged. Even though this technique has shown removing the redundant parameters [LeCun et al., 1990, Hassibi et al., 1993] or the neurons [Mozer and Smolensky, 1988, Ji et al., 1990, Reed, 1993] has little influence on the output, it requires the analysis of each parameter and neuron by network Hessian, which is very computationally expensive for large architectures. The main motivation of these works was to build a more compact network.<br />
<br />
Particularly, building a compact network is a research focus for '''Convolutional Neural Networks'''(CNNs). Some works has proposed to decompose the filters of a pre-trained network into low-rank filters, which reduces the number of parameters [Jaderberg et al., 2014b, Denton et al., 2014, Gong et al., 2014]. The issue of this proposal is that we need to successfully train an initial deep network, since it acts as as post-processing step. [Weigend et al., 1991] and [Collins and Kohl, 2014] used direct training to develop regularizers that eliminate some of the parameters of the network. The problem is that the number of layers and neurons each layer is determined manually. A very similar work using the group lasso method for CNN was previously done in [Liu et al., 2015]. The big-picture idea appears to be very similar but they differ in details of methodology.<br />
<br />
='''Model Training and Model Selection'''=<br />
<br />
In general, a deep network has L layers containing linear operations on their inputs, intertwined with activation functions. The activation function we generally use is '''Rectified Linear Units(RELU) or sigmoids'''. Suppose each layer l has $N_{l}$ neurons, and each of them has parameters $\Theta=(\theta_{l})_{1\leqslant{l}\leqslant{L}}$, where $\theta_{l}=({\theta^n _{l}})_{1\leqslant{n}\leqslant{N_{l}}}$ and where $\theta^n _{l}=[w_{l}^{n},b_{l}^{n}]$ is consisted of linear component $w_{l}^{n}$ and linear bias $b_{l}^{n}$. Given an input $x$, under the linear, on-linear and pooling operations, we obtain the output $\hat{y}=f(x,\theta)$, where $f(*)$ encodes the succession of linear, non-linear and pooling operations.<br />
<br />
At the step of training, we have N input-output pairs ${(x_{i},y_{i})}_{1\leqslant{i}\leqslant{N}}$, and the loss function is given by $\ell(y_{i},f(x_{i},\Theta))$, which compares the predicted output with the ground-truth output. Generally, we choose logistic loss for classification and the square loss for regression. Therefore, learning the parameters of the network is equivalent to solving the optimization of the following:<br />
$$\displaystyle \min_{\Theta}\frac{1}{N}\sum_{i=1}^{N}\ell(y_{i},f(x_{i},\Theta))+\gamma(\Theta),$$ where $\gamma(\Theta)$ represents a regularizer on the network parameters. Our choice for the regularizer can be $\ell_{2}$-norm(i.e, weight decay) or $\ell_{1}$-norm. $\ell_{2}$-norm usually favours small parameter values, and $\ell_{1}$-norm can only delete those irrelevant parameters, but not the neurons. The goal in this paper is to automatically determine the number of neurons of each layer, but neither of the above techniques achieve this goal. Here, we make use of the '''group sparsity''' [Yuan and Lin., 2007] (starting from an overcomplete network and canceling the influence of some neurons). The regularizer, therefore, can be written as $$\gamma(\Theta)=\sum_{l=1}^{L}\beta_{l}\sqrt{P_{l}}\sum_{n=1}^{N_{l}}||\theta_{l}^{n}||_{2},$$ where $P_{l}$ means the size of the vector that includes the parameters of each neuron in layer $l$, and $\beta_{l}$ balances the influence of the penalty. In practice, we found the most effective way to select $\beta$ is a relatively small one for the first few layers, and a larger weight for the remaining layers. The reason we choose a small weight is that it can prevent deleting too much neurons in the first few layers, so that we have enough information for learning the remaining parameters. The original premise of this paper seemed to suggest a new method that was different from both the constructive and destructive methods described above. However, this approach of starting with an overcomplete network and training with group sparsity appears to be no different from destructive methods. The main contribution here is then the regularization function to act on entire neurons, which is in fairness an interesting approach.<br />
<br />
The group sparsity helps us effectively remove some of the neurons, and also standard regularizers on the individual parameters are effective for the generalization purpose [Bartlett, 19996, Krogh and Hertz, 1992, Theodoridis, 2015, Collins and Kohli, 2014]. By this idea, we introduce '''sparse group Lasso''', which considers a more generalised penalty that merges L1 norm in Lasso with the group lasso (i.e. "two-norm"). This leads to the production of a penalty which specifies solutions that are sparse enough both at an individual and group feature levels [1]. It specifies that the regularizer can be written as $$\gamma(\Theta)=\sum_{l=1}^{L}((1-\alpha)\beta_{l}\sqrt{P_{l}}\sum_{n=1}^{N_{l}}||\theta_{l}^{n}||_{2}+\alpha\beta_{l}||\theta_{l}||_{1},$$ where $\alpha\in[0,1]$. We find that if $\alpha=0$, then we have the group sparsity regularizer. In practice, we use both $\alpha=0$ and $\alpha=0.5$ in the experiments.<br />
<br />
This reminds me of the relationships among Lasso regression, Ridge regression and Elastic Net regression. In lasso regression, the penalized residual sum of squares is composed of the regular residual sum of squared plus a L1 regularizer. In ridge regression, its penalized residual sum of squares is composed of the regular residual sum of squared plus a L2 regularizer. Finally, an elastic net regression is a combination of lasso regularizer and ridge regularizer, where its objective function is to optimize parameters by including both L1 and L2 norms. <br />
<br />
To find the optimization, in this paper we use proximal gradient descent [Parikh and Boyed, 2014]. This approach iteratively takes a gradient step of size t with respect to the loss. The following is the algorithm for it: <br />
<br />
We define proximal operator of f as $$prox_{f}(v)=\displaystyle \min_{x}(\frac{1}{2t}||x-v||_{2}^{2}+f(x))$$ <br />
<br />
<br />
Suppose we want to minimize $f(x)+g(x)$, and the proximal gradient method is given by $$x^{(k+1)}=prox_{t^{k}g}(x^{k}-t^{k}\nabla{f}(x^{k})), k=1,2,3...$$ <br />
<br />
Therefore, we can update our parameter by the above method as $$\tilde{\theta}_{l}^{n}=\displaystyle \min_{\theta_{l}^{n}}\frac{1}{2t}||\theta_{l}^{n}-\hat{\theta}_{l}^{n}||_{2}^{2}+\gamma(\Theta),$$<br />
where $\hat{\theta}_{l}^{n}$ is the solution obtained from the general loss gradient. By the derivative of [Simon et al., 2013], we have a closed-form solution for this problem: <br />
$$\tilde{\theta}_{l}^{n}=(1-\frac{t(1-\alpha)\beta_{l}\sqrt{P_{l}}}{||S(\hat{\theta}_{l}^{n},t\alpha\beta_{l})||_{2})})_{+}S(\hat{\theta}_{l}^{n},t\alpha\beta_{l}),$$<br />
where + refers to taking the maximum between the argument and 0, and $S(*)$ is $$S(a,b)=sign(a)(|a|-b)_{+}$$<br />
In practice, we use stochastic gradient descent and work with mini-batches, and then update the variables of all the groups according to the closed-form of $\tilde{\theta}_{l}^{n}$. When the learning steps terminate, we remove the neurons whose parameters have gone to zero. Additionally, when examining fully-connected layers, the neurons acting on output of zeroed-out neurons of the previous layer also become useless, and are removed accordingly.<br />
<br />
='''Experiment'''=<br />
<br />
==='''Set Up'''===<br />
<br />
They use two large-scale image classification datasets, '''ImageNet''' [Russakovsky et al., 2015] and '''Places2-401''' [Zhou et al., 2015]. They also conducted additional experiments on the '''ICDAR''' character recognition dataset of [Jaderberg et al., 2014a]. <br />
<br />
For ImageNet, they used the subset which contains 1000 categories, with 1.2 million training images and 50000 validation images. For Places2-401, it has more than 10 million images with 401 unique scene categories. 5000 to 30000 images are comprised into per category. Both architectures of these two datasets are based on the VGG-B network(BNet) [Simonyan and Zisserman, 2014] and on DecomposeMe8($Dec_{8}$) [ALvarez and Petersson, 2016]. BNet has 10 convolutional layers followed by 3 fully-connected layers. In the experiment, they remove the first 2 fully-connected layers, which we call $BNet^{C}$. $Dec_{8}$ contains 16 convolutional layers with 1D kernels, which can model 8 2D convolutional layers. Both models were trained for a total of 55 epochs with 12000 batches per epoch and a batch size of 48 and 180 for BNet and $Dec_{8}$, respectively. The learning rate was initialized by 0.01 and then multiplied by 0.1. They set $\beta_{l}$=0.102 for the first three layers and $\beta_{l}$=0.255 for the remaining ones.<br />
<br />
For ICDAR dataset, it consists of 185639 training and 5198 test data split into 36 categories. The architecture here starts 6 1D convolutional layers with max-pooling, rather than 3 convolutional layers with a maxout layer [Goodfellow et al., 2013] after each convolution, followed by one fully-connected layer. They call their architecture as Dec3. The model was trained for a total of 45 epochs with a batch size of 256 and 1000 iterations per epoch. The learning rate was initialized by 0.1 and multiplied by 0.1 in the second, seventh and fifteenth epochs. They set $\beta_{l}$=5.1 for the first layer and $\beta_{l}$=10.2 for the remaining ones.<br />
<br />
==='''Results'''===<br />
<br />
[[File:imageNet.png]]<br />
<br />
The above table show the accuracy comparisons between the original architectures and ours. For $Dec_{8}$ on the ImageNet dataset, we evaluated two additional models: $Dec_{8}-640$ with 640 neurons per layer and $Dec_{8}-768$ with 768 neurons per layer. $Dec_{8}-640_{SGL}$ means the sparse group Lasso regularizer with $\alpha=0.5$ and $Dec_{8}-640_{GS}$ represents the group sparsity regularizer. Note that all our architectures yield an improvement over the original network except $Dec_{8}-768$. For instance, Ours-$Bnet_{GS}^{C}$ increases the performance of 1.6% compared to $BNet^{C}$. <br />
<br />
[[File:44.png]]<br />
<br />
[[File:2.png]]<br />
<br />
The above figures reports the reduced percentage of neurons/parameters with our approach for $BNet^{C}$ and $Dec_{8}$. For example, in the first figure, our approach reduces the number of neurons by over 12% and the number of parameters by around 14%, while improving the generalization ability of 1.6%(as indicated by accuracy gap). The left image in the first figure also shows that reduction in number of neurons is spread all the layers with the largest difference in the L10. For $Dec_{8}$, in the second figure, we find when we increase the number of neurons in each layer, the benefits of our approach become more significant. For instance, $Dec_{8}-640$ with group sparsity regularizer reduces the number of neurons by 10%, and of parameters by 12.48%. The left image in the second figure also shows that reduction in number of neurons is spread all the layers. <br />
<br />
[[File:ICDA.png]]<br />
<br />
Finally, the above figure indicates the experiment results for ICDAR dataset. Here, we used the $Dec_{3}$ architecture, where the last two layers initially contain 512 neurons. The accuracy rate for $MaxPllo_{2Dneurons}$ is 83.8%, and accuracy rate for $Dec_{3}$ is 89.3%, which means 1D filters perform better than a network with 2D kernels. Our model on this dataset reduces 38.64% of neurons and totally up to 80% of the number of parameters with group sparsity regularizer.<br />
<br />
All the above results evidence that our algorithm effectively removes the number of parameters and increases the model accuracy. Our algorithm of automatic model selection effectively performs on the classification task.<br />
<br />
='''Analysis on Testing'''=<br />
<br />
Our algorithm does not remove neurons during the training time, however, we remove those neurons after training, which yields a smaller network at test time. This improvement not only reduces the number of parameters of the network, but also decreases the computational memory cost and increases the speed. <br />
<br />
[[File:table2.png]]<br />
<br />
The above table reports the runtime, memory, as well as the percentage of reduced parameters after removing the zeroed-out neurons. The BNet and $Dec_{8}$ were tested on the dataset of ImageNet, while $Dec_{3-GS}$ was tested on the dataset of ICDAR. From the table, we find that all the models for the ImageNet and ICDAR have speeded up the runtime, for example, $Dec_{8}-768_{GS}$ on ImageNet data speeds up the runtime nearly 16% at the batch size of 8, and $Dec_{3}$ on ICDAR data speeds up nearly 50% at natch size of 16. For the percentage of parameters reduced, we find BNet, $Dec_{8}-640_{GS}$ and $Dec_{8}-768_{GS}$ have reduced 12.06%, 26.51%, and 46.73% respectively. More significantly, for $Dec_{3-GS}$, it reduces 82.35% of the parameters. All of these changes show the benefits at the testing time.<br />
<br />
='''Conclusion'''=<br />
<br />
In this paper, we have introduced an approach that relies on group sparsity regularizer. This approach automatically determines the number of neurons in each layer of a deep network. From the experiments, we found our approach not only reduces the number of parameters in our model, but also saves the computation memory and increases the speed at test time. However, the limitation of our approach is that the number of layers in the network remains fixed.<br />
<br />
='''Critique'''=<br />
The authors of the paper state that "...we assume that the parameters of each neuron in layer $l$ are grouped in a vector of size $P_{l}$ and where $\lambda_{l}$ sets the influence of the penalty. Note that, in the general case, this weight can be different for each layer $l$. In practice, however, we found most effective to have<br />
two different weights: a relatively small one for the first few layers, and a larger weight for the<br />
remaining ones. This effectively prevents killing too many neurons in the first few layers, and thus<br />
retains enough information for the remaining ones." However, the authors fail to present any guidance as to what gets counted as "the first few layers" and what the relative sizes for the two weights should be even after we have chosen the "first few layers". Indeed, such choice seems to be an unaccounted component of tuning the model but this receives scant attention in the current paper. Several numerical comparisons should be carried out to allow further discussion on this question.<br />
<br />
The experiments could have included better baseline models to compare against. For example, how do we know the original model was not overly complex to begin with? It might have been a good idea for the authors to compare their group sparse lasso method against the naive method of (blindly) reducing the number of neurons in each layer by 10-20% just for a very preliminary check. Also, the author could have applied the theory of ridge and Lasso regression to analyze the effect of the regularization mathematically.<br />
<br />
A rather reliable method of experimentation to compare the performance and accuracy has been left out. The authors have not stated any comparisons of this method with the Dropout method [Srivastava,2014], which is similar in terms of the physical effects on the network.<br />
<br />
='''References'''=<br />
<br />
P. L. Bartlett. For valid generalization the size of the weights is more important than the size of the network. In NIPS, 1996.<br />
<br />
M. G. Bello. Enhanced training algorithms, and integrated training/architecture selection for multilayer perceptron networks. IEEE Transactions on Neural Networks, 3(6):864–875, Nov 1992.<br />
<br />
Yu Cheng, Felix X. Yu, Rogério Schmidt Feris, Sanjiv Kumar, Alok N. Choudhary, and Shih-Fu Chang. An exploration of parameter redundancy in deep networks with circulant projections. In ICCV, 2015.<br />
<br />
I. J. Goodfellow, D. Warde-farley, M. Mirza, A. Courville, and Y. Bengio. Maxout networks. In ICML, 2013.<br />
<br />
G. E. Hinton, O. Vinyals, and J. Dean. Distilling the knowledge in a neural network. In arXiv, 2014.<br />
<br />
M. Jaderberg, A. Vedaldi, and A. Zisserman. Deep features for text spotting. In ECCV, 2014a.<br />
<br />
M. Jaderberg, A. Vedaldi, and A. Zisserman. Speeding up convolutional neural networks with low rank expansions. In BMVC, 2014b.<br />
<br />
N. Simon, J. Friedman, T. Hastie, and R. Tibshirani. A sparse-group lasso. Journal of Computational and Graphical Statistics, 2013.<br />
<br />
H. Zhou, J. M. Alvarez, and F. Porikli. Less is more: Towards compact CNNs. In ECCV, 2016.<br />
<br />
Group LASSO - https://pdfs.semanticscholar.org/f677/a011b2a912e3c5c604f6872b9716cc0b8aa0.pdf<br />
<br />
Liu, Baoyuan, et al. "Sparse convolutional neural networks." Proceedings of the IEEE Conference on Computer Vision and Pattern Recognition. 2015.<br />
<br />
Nitish Srivastava, Geoffrey Hinton, Alex Krizhevsky, Ilya Sutskever, and Ruslan Salakhutdinov. 2014. Dropout: a simple way to prevent neural networks from overfitting. J. Mach. Learn. Res. 15, 1 (January 2014), 1929-1958.<br />
<br />
<br />
Derivation & Motivation of the Soft Thresholding Operator (Proximal Operator):<br />
# http://www.onmyphd.com/?p=proximal.operator<br />
# https://math.stackexchange.com/questions/471339/derivation-of-soft-thresholding-operator</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/Cognitive_Psychology_For_Deep_Neural_Networks:_A_Shape_Bias_Case_Study&diff=29765STAT946F17/Cognitive Psychology For Deep Neural Networks: A Shape Bias Case Study2017-11-09T00:57:45Z<p>Nkadiyar: /* Cognitive Biases */</p>
<hr />
<div>= Introduction =<br />
<br />
The recent burgeon on the use of Deep Neural Networks (DNNs) have resulted in giant leaps of accuracy in prediction. They are also being used to solve a variety of complex tasks which earlier methodologies have struggled to excel in.<br />
<br />
While it is all good to see incredibly high accuracy as a result of the use of DNN, we must begin to question why they perform so well. It has become an interesting field of study to actually represent the features/feature maps or interpret the meaning of the learnt values in a DNN's hidden layers. Currently we treat models of DNNs as black boxes which we practically tune the tweakable parameters like number of layers, number of units in each layer, number & size of feature maps(in case of CNN) etc. The opacity created by the lack of an intuitive representation of the internal learnt parameters of DNNs hinders both basic research as well as its application to real world problems.<br />
<br />
Recent pushes have aimed to better understand DNNs: tailor-made loss functions and architectures produce more interpretable features (Higgins et al., 2016; Raposo et al., 2017) while output-behavior analyses unveil previously opaque operations of these networks (Karpathy et al., 2015). Parallel to this work, neuroscience-inspired methods such as activation visualization (Li et al., 2015), ablation analysis (Zeiler & Fergus, 2014) and activation maximization (Yosinski et al., 2015) have also been applied<br />
<br />
This paper aims to provide another methodology to attempt to decipher & better understand how DNNs solve a particular task. This methodology was inspired by psychological concepts to test whether the DNN's were able to make accurate predictions with biases similar to that the human mind makes.<br />
<br />
Research in developmental psychology shows that when learning new words, humans tend to assign the same name to similarly shaped items rather than to items with similar color, texture, or size. This bias/knowledge tend to be forged into the brains of humans and humans then take this forward to easily associate these shapes with new objects they have not seen before.<br />
<br />
The authors of this paper try to test if DNNs behave similarly in one-shot learning applications. They attempt to prove that when the models of state-of-the-art DNNs are used to learn objects from images, they exhibit a stronger shape bias than a color bias. To emulate the human brain, they use the parameters of pre-trained DNN models and use this to perform one-shot learning on a new data set with different labels.<br />
<br />
= Background =<br />
== One Shot Learning ==<br />
One-shot learning is an object categorization problem in computer vision. Whereas most machine learning based object categorization algorithms require training on hundreds or thousands of images and very large datasets, one-shot learning aims to learn information about object categories from one, or only a few, training images.<br />
<br />
The one-shot word learning task is to label a novel data example $\hat{x}$ (e.g. a novel probe image) with a novel class label $\hat{y}$ (e.g. a new word) after only a single example.<br />
<br />
More specifically, given a support set $S = {(x_i, y_i) , i \in [1, k]}$, of images $x_i$, and their associated labels $y_i$, and an unlabeled probe image $\hat{x}$,<br />
the one-shot learning task is to identify the true label of the probe image, $\hat{y}$, from the support set labels $ {y_i , i \in [1, k]} $:<br />
<br />
<br />
$\displaystyle \hat{y} = arg \max_{y}$ $P(y | \hat{x}, S)$<br />
<br />
<br />
We assume that the image labels $y_i$ are represented using a one-hot encoding and that $P(y|\hat{x}, S)$ is parameterised by a DNN, allowing us to leverage the ability of deep networks to learn powerful representations.<br />
<br />
== Inception Networks ==<br />
<br />
A probe image $\hat{x}$ is given the label of the nearest neighbour from the<br />
support set:<br />
<br />
$\hat{y} = y$<br />
<br />
$(x, y) = \displaystyle arg \min_{(x_i,y_i) \in S} d(h(x_i), h(\hat{x})) $<br />
<br />
where d is a distance function.<br />
<br />
The function h is parameterized by Inception – one of the best performing ImageNet classification models. Specifically, h returns features from the last layer (the softmax input) of a pre-trained Inception classifier. With these features as input and cosine distance as the distance function, the classifier in achieves 87.6% accuracy on one-shot classification on the ImageNet dataset (Vinyals et al., 2016). We call the Inception classifier together with the nearest-neighbor component the Inception Baseline (IB) model.<br />
<br />
<br />
== Matching Networks ==<br />
<br />
MNs (Vinyals et al.,2016) are neural network architectures with state-of-the-art one shot learning performance on ImageNet (93.2% one-shot labelling accuracy).<br />
MNs are trained to assign label $\hat{y}$ to probe image $\hat{x}$ using an attention mechanism $a$ acting on image embeddings stored in the support set S:<br />
\begin{align*}<br />
a(\hat{x},x_i)=\frac{e^{d(f(\hat{x},S), g(x_i,S))}}{\sum_{j}e^{d(f(\hat{x},S), g(x_j,S))}},<br />
\end{align*}<br />
<br />
<br />
where d is a cosine distance and where f and g provide context-dependent embeddings of $\hat{x}$ and $x_i$ (with contextS). The embedding $g(x_i, S)$ is a bi-directional LSTM (Hochreiter & Schmidhuber, 1997) with the support set S provided as an input sequence. The embedding $f(\hat{x}, S)$ is an LSTM with a read-attention mechanism operating over the entire embedded support set. The input to the LSTM is given by the penultimate layer features of a pre-trained deep convolutional network, specifically Inception.<br />
<br />
To train MNs we proceed as follows:<br />
<br />
=== Training MN ===<br />
* Step 1: At each step of training, the model is given a small support set of images and associated labels. In addition to the support set, the model is fed an unlabeled probe image $\hat{x}$<br />
<br />
* Step 2: The model parameters are then updated to improve classification accuracy of the probe image $\hat{x}$ given the support set. Parameters are updated using stochastic gradient descent with a learning rate of 0.1<br />
<br />
* Step 3: After each update, the labels ${(y_i, i \in [1, k]}$ in the training set are randomly re-assigned to new image classes (the label indices are randomly permuted,<br />
but the image labels are not changed). This is a critical step. It prevents MNs from learning a consistent mapping between a category and a label. Usually, in classification, this is what we want, but in one-shot learning we want to train our model for classification after viewing a single in-class example from the support set. <br />
<br />
The objective function used is:<br />
<br />
\begin{align*}<br />
L=E_{C\sim T}\biggr[E_{S \sim C, B \sim C}\bigr[\sum_{(x,y)\in B}\log P(y|x,S)\bigr]\biggr]<br />
\end{align*}<br />
<br />
<br />
where T is the set of all possible labelings of our classes, S is a support set sampled with a class labeling C ~ T and B is a batch of probe images and labels, also with the same randomly chosen class labeling as the support set.<br />
<br />
== Cognitive Biases ==<br />
A Cognitive bias, in simple terms, is any systematic deviation from logic. Wikipedia defines a Cognitive bias as "a systematic pattern of deviation from norm or rationality in judgment, whereby inferences about other people and situations may be drawn in an illogical fashion". They can be seen as mental shortcuts (although not always correct) that the mind uses to make decisions or learn. Cognitive bias is a concept from developmental psychology which attempts to explain how children can extract meanings of words with very few examples, similar to the concept of one-shot learning discussed above. The theory, as explained by the authors, is that humans form biases that allow them to eliminate many potential hypotheses about word meaning where the amount of data available is insufficient for this purpose. These include:<br />
* Whole object bias<br />
* Taxonomic bias<br />
* Mutual exclusivity bias<br />
* Shape bias<br />
A more complete list of cognitive biases is given by [[#References|(Bloom, 2000)]]. The bias the authors investigate in this paper is the shape bias, which denotes a tendency to assign the same name to similarly shaped items rather than to items with similar color, texture, or size.<br />
<br />
= Methodology =<br />
== Inductive Biases & Probe Data ==<br />
<br />
Inductive biases are those criteria which are artificially selected or learnt by the network as a classifying/distinguishing property.<br />
It has been observed that the biases that DNNs learnt are complex composite features. We, as researchers can take advantage of the fact that DNNs learnt complex distinguishing features by constructing probe data sets which particularly target on exposing a particular bias that a DNN might have. <br />
<br />
* Step 1: Take a known composite feature which we suspect the DNNs are biased against. It is possible that the feature is not numerical but intuitive for human researchers to understand.<br />
* Step 2: Train the target model with an appropriate dataset.<br />
* Step 3: Transfer Learning: Use the pre-trained model with a new data set which is curated to contain data to prove/disprove the existence of the bias<br />
* Step 4: Model/Decide on a function which quantifies the bias under study<br />
* Step 5: Measure the bias with the bias function<br />
<br />
Then it is possible to deduce whether the DNN uses the feature to solve the task by the value of bias function.<br />
<br />
== Data Sets Used ==<br />
<br />
* Training Set: ImageNet<br />
* Test Set:<br />
** The Cognitive Psychology Probe Data (CogPsyc data) that is used consists of 150 images of objects. The images are arranged in triples consisting of a probe image, a shape-match image (that matches the probe in colour but not shape), and a color-match image (that matches the probe in shape but not colour). In the dataset there are 10 triples, each shown on 5 different backgrounds, giving a total of 50 triples. [[File:CogPsy.PNG|center|350px]]<br />
** A real-world dataset consisting of 90 images of objects (30 triples) collected using Google Image Search. The images are arranged in triples consisting of a probe, a shape-match and a colour-match.<br />
<br />
= Experiments =<br />
== Evaluation Criteria ==<br />
<br />
* For a given probe image $\hat{x}$, we loaded the shape-match image $x_s$ and corresponding label $y_s$, along with the colour-match image $x_c$ and corresponding label $y_c$ into memory, as the support set $S = \{(x_s, ys), (x_c, y_c)\}$<br />
* Calculate $\hat{y}$<br />
* The model assigns either $y_c$ or $y_s$ to the probe image.<br />
* To estimate the shape bias Bs, calculate the proportion of shape labels assigned to the probe: $B_s = E(\delta(\hat{y} - y_s))$ <br />
where E is an expectation across probe images and $\delta$ is the Dirac delta function.<br />
<br />
== Experiment 1: Shape bias statistics in Inception Baseline: ==<br />
* Shape bias of IB to be $B_s = 0.68$. Similarly, the shape bias of IB using our real-world dataset was $B_s = 0.97$. Together, these results strongly suggest that IB trained on ImageNet has a stronger bias towards shape than colour<br />
<br />
== Experiment 2: Shape bias statistics in Matching Network: ==<br />
* They found that MNs have a shape of bias $B_s = 0.7$ using the CogPsyc dataset and a bias of $Bs = 1$ using the real-world dataset. Once again, these results suggest that MNs trained seeding from Inception using ImageNet has a stronger bias towards shape than colour.<br />
<br />
== Experiment 3: Shape bias statistics between and across models: ==<br />
<br />
The authors extended the shape bias analysis to calculate the shape bias in a population of IB models and in a population of MN models with different random initialization<br />
<br />
=== Dependence on the initialization of parameters: ===<br />
<br />
[[File:3.1.PNG|right|250px]] A strong variability was observed when variation in the initial values of the parameters. For the CogPsyc dataset, the average shape bias was $B_s = 0.628$ with standard deviation $\sigma B_s = 0.049$ at the end of training and for the real-world dataset the average shape bias was $B_s = 0:958$ with $\sigma B_s = 0.037$.<br />
<br />
=== Dependence of shape bias on model performance: ===<br />
<br />
For the CogPsych dataset, the correlation between bias and classification accuracy was $\rho = 0.15$, and for the real world dataset, correlation between bias and classification accuracy was $\rho = -0.06$. This would be evident since the accuracy of the models remained nearly constant when the initialization parameters varied whereas the shape bias tended to vary a lot, hence highlighting the lack of correlation amongst them. <br />
<br />
=== Emergence of shape bias during training: ===<br />
The shape bias spiked to a large value very early. <br />
<br />
=== Variation of shape bias within models & across models: ===<br />
With different initialization parameters, the shape bias varied a lot within IB during training while the shape bias did not fluctuate during the training of MN. It was found that the MN inherits the shape bias of the IB which seeded its embeddings and thereafter, the shape bias remained constant throughout training. It is important to note that the output of the penultimate layer of the Inception was not fine tuned when it was pipelined to the MN. This was to ensure that the MN properties were independent of the IB model properties. [[File:3.3.PNG|center|250px]] [[File:3.4.PNG|center|250px]]<br />
<br />
= Learnings, Inferences & Implications =<br />
* Both the Inception Baseline and the Matching Network exhibit strong shape bias when trained on ImageNet. Researchers who use Inception & MN DNNs can now use this fact as a consideration for their application while using pre-trained models for new datasets. If it is known before hand that the new data set is strongly classifiable through a color bias, then they would either want to defer using the pre-trained models or explore methods to decrease/remove the strong shape bias.<br />
<br />
* There exists a high variability in the shape bias with the variation in the initialization parameters. This is an important finding since it uncovers the fact that the same architecture which exhibit similar accuracy in predictions can display a variety of shape bias just with different initialization parameters. Researchers can explore methods of tuning the random initialization such that the models start out with a low shape bias without compromising the accuracy of the model.<br />
<br />
* MNs inherit the shape bias which is seeded to it by the Inception Network's input embedding. This is also another fact which researchers & practitioners should be careful about. When using cascaded or pipelined heterogeneous architectures, the models downstream tend to inherit/become/are fed with the properties/biases of the models upstream. This may be desirable or undesirable according to the application, but it is important to be aware of its presence.<br />
<br />
* The biases under consideration are the property of the collection of the architecture, the dataset and the optimization procedure. Hence in order to increase or decrease the effect of a particular bias, one or more of the mentioned factors must be adjusted/tuned/changed.<br />
<br />
* The fact that a high shape bias emerged in the early epochs with less variability in further epochs can be thought of analogous to the biases that humans develop at an infancy which gets fortified as they age.<br />
<br />
=== Modeling human word learning === <br />
The authors note how there have been several previous attempts to model human word learning in the field of cognitive science. A major shortcoming of these works is that none of the models are capable of one-shot word learning when presented with real-world images. Recognizing this and given the success of MNs in this work, the authors propose MNs as a computational-level account of one-shot word learning. <br />
In the paper, it is discussed how shape bias increases dramatically during early training of the model. This agrees with what is observed in many psychological studies: older children show bias more than younger children. <br />
<br />
= Conclusion, Future Work and Open questions =<br />
<br />
* Just as cognitive psychology exposes the shape bias observed in this experiment, we should try to uncover other biases as well using multiple approaches<br />
* Study the underlying mechanisms which cause biases such as shape bias in DNNs<br />
* Research into various methods of probing and creating probe data sets which can be used to test architectures for various biases<br />
* Exploration into a research field called Artificial Cognitive Psychology which focuses on probing how DNN architectures can be understood further using known behaviors of the human brain<br />
<br />
= References =<br />
<br />
* Ritter, Samuel & G. T. Barrett, David & Santoro, Adam & M. Botvinick, Matt. (2017). Cognitive Psychology for Deep Neural Networks: A Shape Bias Case Study<br />
<br />
* Vinyals, Oriol, Blundell, Charles, Lillicrap, Timothy, Kavukcuoglu, Koray, and Wierstra, Daan. Matching networks for one shot learning. arXiv preprint arXiv:1606.04080, 2016.<br />
<br />
* Bloom, P. (2000). How children learn the meanings of words. The MIT Press.<br />
<br />
* https://www.slideshare.net/KazukiFujikawa/matching-networks-for-one-shot-learning-71257100<br />
<br />
* https://deepmind.com/blog/cognitive-psychology/<br />
<br />
* https://hacktilldawn.com/2016/09/25/inception-modules-explained-and-implemented/</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/Cognitive_Psychology_For_Deep_Neural_Networks:_A_Shape_Bias_Case_Study&diff=29764STAT946F17/Cognitive Psychology For Deep Neural Networks: A Shape Bias Case Study2017-11-09T00:56:56Z<p>Nkadiyar: /* Cognitive Biases */</p>
<hr />
<div>= Introduction =<br />
<br />
The recent burgeon on the use of Deep Neural Networks (DNNs) have resulted in giant leaps of accuracy in prediction. They are also being used to solve a variety of complex tasks which earlier methodologies have struggled to excel in.<br />
<br />
While it is all good to see incredibly high accuracy as a result of the use of DNN, we must begin to question why they perform so well. It has become an interesting field of study to actually represent the features/feature maps or interpret the meaning of the learnt values in a DNN's hidden layers. Currently we treat models of DNNs as black boxes which we practically tune the tweakable parameters like number of layers, number of units in each layer, number & size of feature maps(in case of CNN) etc. The opacity created by the lack of an intuitive representation of the internal learnt parameters of DNNs hinders both basic research as well as its application to real world problems.<br />
<br />
Recent pushes have aimed to better understand DNNs: tailor-made loss functions and architectures produce more interpretable features (Higgins et al., 2016; Raposo et al., 2017) while output-behavior analyses unveil previously opaque operations of these networks (Karpathy et al., 2015). Parallel to this work, neuroscience-inspired methods such as activation visualization (Li et al., 2015), ablation analysis (Zeiler & Fergus, 2014) and activation maximization (Yosinski et al., 2015) have also been applied<br />
<br />
This paper aims to provide another methodology to attempt to decipher & better understand how DNNs solve a particular task. This methodology was inspired by psychological concepts to test whether the DNN's were able to make accurate predictions with biases similar to that the human mind makes.<br />
<br />
Research in developmental psychology shows that when learning new words, humans tend to assign the same name to similarly shaped items rather than to items with similar color, texture, or size. This bias/knowledge tend to be forged into the brains of humans and humans then take this forward to easily associate these shapes with new objects they have not seen before.<br />
<br />
The authors of this paper try to test if DNNs behave similarly in one-shot learning applications. They attempt to prove that when the models of state-of-the-art DNNs are used to learn objects from images, they exhibit a stronger shape bias than a color bias. To emulate the human brain, they use the parameters of pre-trained DNN models and use this to perform one-shot learning on a new data set with different labels.<br />
<br />
= Background =<br />
== One Shot Learning ==<br />
One-shot learning is an object categorization problem in computer vision. Whereas most machine learning based object categorization algorithms require training on hundreds or thousands of images and very large datasets, one-shot learning aims to learn information about object categories from one, or only a few, training images.<br />
<br />
The one-shot word learning task is to label a novel data example $\hat{x}$ (e.g. a novel probe image) with a novel class label $\hat{y}$ (e.g. a new word) after only a single example.<br />
<br />
More specifically, given a support set $S = {(x_i, y_i) , i \in [1, k]}$, of images $x_i$, and their associated labels $y_i$, and an unlabeled probe image $\hat{x}$,<br />
the one-shot learning task is to identify the true label of the probe image, $\hat{y}$, from the support set labels $ {y_i , i \in [1, k]} $:<br />
<br />
<br />
$\displaystyle \hat{y} = arg \max_{y}$ $P(y | \hat{x}, S)$<br />
<br />
<br />
We assume that the image labels $y_i$ are represented using a one-hot encoding and that $P(y|\hat{x}, S)$ is parameterised by a DNN, allowing us to leverage the ability of deep networks to learn powerful representations.<br />
<br />
== Inception Networks ==<br />
<br />
A probe image $\hat{x}$ is given the label of the nearest neighbour from the<br />
support set:<br />
<br />
$\hat{y} = y$<br />
<br />
$(x, y) = \displaystyle arg \min_{(x_i,y_i) \in S} d(h(x_i), h(\hat{x})) $<br />
<br />
where d is a distance function.<br />
<br />
The function h is parameterized by Inception – one of the best performing ImageNet classification models. Specifically, h returns features from the last layer (the softmax input) of a pre-trained Inception classifier. With these features as input and cosine distance as the distance function, the classifier in achieves 87.6% accuracy on one-shot classification on the ImageNet dataset (Vinyals et al., 2016). We call the Inception classifier together with the nearest-neighbor component the Inception Baseline (IB) model.<br />
<br />
<br />
== Matching Networks ==<br />
<br />
MNs (Vinyals et al.,2016) are neural network architectures with state-of-the-art one shot learning performance on ImageNet (93.2% one-shot labelling accuracy).<br />
MNs are trained to assign label $\hat{y}$ to probe image $\hat{x}$ using an attention mechanism $a$ acting on image embeddings stored in the support set S:<br />
\begin{align*}<br />
a(\hat{x},x_i)=\frac{e^{d(f(\hat{x},S), g(x_i,S))}}{\sum_{j}e^{d(f(\hat{x},S), g(x_j,S))}},<br />
\end{align*}<br />
<br />
<br />
where d is a cosine distance and where f and g provide context-dependent embeddings of $\hat{x}$ and $x_i$ (with contextS). The embedding $g(x_i, S)$ is a bi-directional LSTM (Hochreiter & Schmidhuber, 1997) with the support set S provided as an input sequence. The embedding $f(\hat{x}, S)$ is an LSTM with a read-attention mechanism operating over the entire embedded support set. The input to the LSTM is given by the penultimate layer features of a pre-trained deep convolutional network, specifically Inception.<br />
<br />
To train MNs we proceed as follows:<br />
<br />
=== Training MN ===<br />
* Step 1: At each step of training, the model is given a small support set of images and associated labels. In addition to the support set, the model is fed an unlabeled probe image $\hat{x}$<br />
<br />
* Step 2: The model parameters are then updated to improve classification accuracy of the probe image $\hat{x}$ given the support set. Parameters are updated using stochastic gradient descent with a learning rate of 0.1<br />
<br />
* Step 3: After each update, the labels ${(y_i, i \in [1, k]}$ in the training set are randomly re-assigned to new image classes (the label indices are randomly permuted,<br />
but the image labels are not changed). This is a critical step. It prevents MNs from learning a consistent mapping between a category and a label. Usually, in classification, this is what we want, but in one-shot learning we want to train our model for classification after viewing a single in-class example from the support set. <br />
<br />
The objective function used is:<br />
<br />
\begin{align*}<br />
L=E_{C\sim T}\biggr[E_{S \sim C, B \sim C}\bigr[\sum_{(x,y)\in B}\log P(y|x,S)\bigr]\biggr]<br />
\end{align*}<br />
<br />
<br />
where T is the set of all possible labelings of our classes, S is a support set sampled with a class labeling C ~ T and B is a batch of probe images and labels, also with the same randomly chosen class labeling as the support set.<br />
<br />
== Cognitive Biases ==<br />
A Cognitive bias, in simple terms, is any systematic deviation from logic. Wikipedia defines a Cognitive bias as "a systematic pattern of deviation from norm or rationality in judgment, whereby inferences about other people and situations may be drawn in an illogical fashion". They can be seen as mental shortcuts (although not always correct) that the mind uses to judge or learn. Cognitive bias is a concept from developmental psychology which attempts to explain how children can extract meanings of words with very few examples, similar to the concept of one-shot learning discussed above. The theory, as explained by the authors, is that humans form biases that allow them to eliminate many potential hypotheses about word meaning where the amount of data available is insufficient for this purpose. These include:<br />
* Whole object bias<br />
* Taxonomic bias<br />
* Mutual exclusivity bias<br />
* Shape bias<br />
A more complete list of cognitive biases is given by [[#References|(Bloom, 2000)]]. The bias the authors investigate in this paper is the shape bias, which denotes a tendency to assign the same name to similarly shaped items rather than to items with similar color, texture, or size.<br />
<br />
= Methodology =<br />
== Inductive Biases & Probe Data ==<br />
<br />
Inductive biases are those criteria which are artificially selected or learnt by the network as a classifying/distinguishing property.<br />
It has been observed that the biases that DNNs learnt are complex composite features. We, as researchers can take advantage of the fact that DNNs learnt complex distinguishing features by constructing probe data sets which particularly target on exposing a particular bias that a DNN might have. <br />
<br />
* Step 1: Take a known composite feature which we suspect the DNNs are biased against. It is possible that the feature is not numerical but intuitive for human researchers to understand.<br />
* Step 2: Train the target model with an appropriate dataset.<br />
* Step 3: Transfer Learning: Use the pre-trained model with a new data set which is curated to contain data to prove/disprove the existence of the bias<br />
* Step 4: Model/Decide on a function which quantifies the bias under study<br />
* Step 5: Measure the bias with the bias function<br />
<br />
Then it is possible to deduce whether the DNN uses the feature to solve the task by the value of bias function.<br />
<br />
== Data Sets Used ==<br />
<br />
* Training Set: ImageNet<br />
* Test Set:<br />
** The Cognitive Psychology Probe Data (CogPsyc data) that is used consists of 150 images of objects. The images are arranged in triples consisting of a probe image, a shape-match image (that matches the probe in colour but not shape), and a color-match image (that matches the probe in shape but not colour). In the dataset there are 10 triples, each shown on 5 different backgrounds, giving a total of 50 triples. [[File:CogPsy.PNG|center|350px]]<br />
** A real-world dataset consisting of 90 images of objects (30 triples) collected using Google Image Search. The images are arranged in triples consisting of a probe, a shape-match and a colour-match.<br />
<br />
= Experiments =<br />
== Evaluation Criteria ==<br />
<br />
* For a given probe image $\hat{x}$, we loaded the shape-match image $x_s$ and corresponding label $y_s$, along with the colour-match image $x_c$ and corresponding label $y_c$ into memory, as the support set $S = \{(x_s, ys), (x_c, y_c)\}$<br />
* Calculate $\hat{y}$<br />
* The model assigns either $y_c$ or $y_s$ to the probe image.<br />
* To estimate the shape bias Bs, calculate the proportion of shape labels assigned to the probe: $B_s = E(\delta(\hat{y} - y_s))$ <br />
where E is an expectation across probe images and $\delta$ is the Dirac delta function.<br />
<br />
== Experiment 1: Shape bias statistics in Inception Baseline: ==<br />
* Shape bias of IB to be $B_s = 0.68$. Similarly, the shape bias of IB using our real-world dataset was $B_s = 0.97$. Together, these results strongly suggest that IB trained on ImageNet has a stronger bias towards shape than colour<br />
<br />
== Experiment 2: Shape bias statistics in Matching Network: ==<br />
* They found that MNs have a shape of bias $B_s = 0.7$ using the CogPsyc dataset and a bias of $Bs = 1$ using the real-world dataset. Once again, these results suggest that MNs trained seeding from Inception using ImageNet has a stronger bias towards shape than colour.<br />
<br />
== Experiment 3: Shape bias statistics between and across models: ==<br />
<br />
The authors extended the shape bias analysis to calculate the shape bias in a population of IB models and in a population of MN models with different random initialization<br />
<br />
=== Dependence on the initialization of parameters: ===<br />
<br />
[[File:3.1.PNG|right|250px]] A strong variability was observed when variation in the initial values of the parameters. For the CogPsyc dataset, the average shape bias was $B_s = 0.628$ with standard deviation $\sigma B_s = 0.049$ at the end of training and for the real-world dataset the average shape bias was $B_s = 0:958$ with $\sigma B_s = 0.037$.<br />
<br />
=== Dependence of shape bias on model performance: ===<br />
<br />
For the CogPsych dataset, the correlation between bias and classification accuracy was $\rho = 0.15$, and for the real world dataset, correlation between bias and classification accuracy was $\rho = -0.06$. This would be evident since the accuracy of the models remained nearly constant when the initialization parameters varied whereas the shape bias tended to vary a lot, hence highlighting the lack of correlation amongst them. <br />
<br />
=== Emergence of shape bias during training: ===<br />
The shape bias spiked to a large value very early. <br />
<br />
=== Variation of shape bias within models & across models: ===<br />
With different initialization parameters, the shape bias varied a lot within IB during training while the shape bias did not fluctuate during the training of MN. It was found that the MN inherits the shape bias of the IB which seeded its embeddings and thereafter, the shape bias remained constant throughout training. It is important to note that the output of the penultimate layer of the Inception was not fine tuned when it was pipelined to the MN. This was to ensure that the MN properties were independent of the IB model properties. [[File:3.3.PNG|center|250px]] [[File:3.4.PNG|center|250px]]<br />
<br />
= Learnings, Inferences & Implications =<br />
* Both the Inception Baseline and the Matching Network exhibit strong shape bias when trained on ImageNet. Researchers who use Inception & MN DNNs can now use this fact as a consideration for their application while using pre-trained models for new datasets. If it is known before hand that the new data set is strongly classifiable through a color bias, then they would either want to defer using the pre-trained models or explore methods to decrease/remove the strong shape bias.<br />
<br />
* There exists a high variability in the shape bias with the variation in the initialization parameters. This is an important finding since it uncovers the fact that the same architecture which exhibit similar accuracy in predictions can display a variety of shape bias just with different initialization parameters. Researchers can explore methods of tuning the random initialization such that the models start out with a low shape bias without compromising the accuracy of the model.<br />
<br />
* MNs inherit the shape bias which is seeded to it by the Inception Network's input embedding. This is also another fact which researchers & practitioners should be careful about. When using cascaded or pipelined heterogeneous architectures, the models downstream tend to inherit/become/are fed with the properties/biases of the models upstream. This may be desirable or undesirable according to the application, but it is important to be aware of its presence.<br />
<br />
* The biases under consideration are the property of the collection of the architecture, the dataset and the optimization procedure. Hence in order to increase or decrease the effect of a particular bias, one or more of the mentioned factors must be adjusted/tuned/changed.<br />
<br />
* The fact that a high shape bias emerged in the early epochs with less variability in further epochs can be thought of analogous to the biases that humans develop at an infancy which gets fortified as they age.<br />
<br />
=== Modeling human word learning === <br />
The authors note how there have been several previous attempts to model human word learning in the field of cognitive science. A major shortcoming of these works is that none of the models are capable of one-shot word learning when presented with real-world images. Recognizing this and given the success of MNs in this work, the authors propose MNs as a computational-level account of one-shot word learning. <br />
In the paper, it is discussed how shape bias increases dramatically during early training of the model. This agrees with what is observed in many psychological studies: older children show bias more than younger children. <br />
<br />
= Conclusion, Future Work and Open questions =<br />
<br />
* Just as cognitive psychology exposes the shape bias observed in this experiment, we should try to uncover other biases as well using multiple approaches<br />
* Study the underlying mechanisms which cause biases such as shape bias in DNNs<br />
* Research into various methods of probing and creating probe data sets which can be used to test architectures for various biases<br />
* Exploration into a research field called Artificial Cognitive Psychology which focuses on probing how DNN architectures can be understood further using known behaviors of the human brain<br />
<br />
= References =<br />
<br />
* Ritter, Samuel & G. T. Barrett, David & Santoro, Adam & M. Botvinick, Matt. (2017). Cognitive Psychology for Deep Neural Networks: A Shape Bias Case Study<br />
<br />
* Vinyals, Oriol, Blundell, Charles, Lillicrap, Timothy, Kavukcuoglu, Koray, and Wierstra, Daan. Matching networks for one shot learning. arXiv preprint arXiv:1606.04080, 2016.<br />
<br />
* Bloom, P. (2000). How children learn the meanings of words. The MIT Press.<br />
<br />
* https://www.slideshare.net/KazukiFujikawa/matching-networks-for-one-shot-learning-71257100<br />
<br />
* https://deepmind.com/blog/cognitive-psychology/<br />
<br />
* https://hacktilldawn.com/2016/09/25/inception-modules-explained-and-implemented/</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/Cognitive_Psychology_For_Deep_Neural_Networks:_A_Shape_Bias_Case_Study&diff=29763STAT946F17/Cognitive Psychology For Deep Neural Networks: A Shape Bias Case Study2017-11-09T00:53:03Z<p>Nkadiyar: /* Cognitive Biases */</p>
<hr />
<div>= Introduction =<br />
<br />
The recent burgeon on the use of Deep Neural Networks (DNNs) have resulted in giant leaps of accuracy in prediction. They are also being used to solve a variety of complex tasks which earlier methodologies have struggled to excel in.<br />
<br />
While it is all good to see incredibly high accuracy as a result of the use of DNN, we must begin to question why they perform so well. It has become an interesting field of study to actually represent the features/feature maps or interpret the meaning of the learnt values in a DNN's hidden layers. Currently we treat models of DNNs as black boxes which we practically tune the tweakable parameters like number of layers, number of units in each layer, number & size of feature maps(in case of CNN) etc. The opacity created by the lack of an intuitive representation of the internal learnt parameters of DNNs hinders both basic research as well as its application to real world problems.<br />
<br />
Recent pushes have aimed to better understand DNNs: tailor-made loss functions and architectures produce more interpretable features (Higgins et al., 2016; Raposo et al., 2017) while output-behavior analyses unveil previously opaque operations of these networks (Karpathy et al., 2015). Parallel to this work, neuroscience-inspired methods such as activation visualization (Li et al., 2015), ablation analysis (Zeiler & Fergus, 2014) and activation maximization (Yosinski et al., 2015) have also been applied<br />
<br />
This paper aims to provide another methodology to attempt to decipher & better understand how DNNs solve a particular task. This methodology was inspired by psychological concepts to test whether the DNN's were able to make accurate predictions with biases similar to that the human mind makes.<br />
<br />
Research in developmental psychology shows that when learning new words, humans tend to assign the same name to similarly shaped items rather than to items with similar color, texture, or size. This bias/knowledge tend to be forged into the brains of humans and humans then take this forward to easily associate these shapes with new objects they have not seen before.<br />
<br />
The authors of this paper try to test if DNNs behave similarly in one-shot learning applications. They attempt to prove that when the models of state-of-the-art DNNs are used to learn objects from images, they exhibit a stronger shape bias than a color bias. To emulate the human brain, they use the parameters of pre-trained DNN models and use this to perform one-shot learning on a new data set with different labels.<br />
<br />
= Background =<br />
== One Shot Learning ==<br />
One-shot learning is an object categorization problem in computer vision. Whereas most machine learning based object categorization algorithms require training on hundreds or thousands of images and very large datasets, one-shot learning aims to learn information about object categories from one, or only a few, training images.<br />
<br />
The one-shot word learning task is to label a novel data example $\hat{x}$ (e.g. a novel probe image) with a novel class label $\hat{y}$ (e.g. a new word) after only a single example.<br />
<br />
More specifically, given a support set $S = {(x_i, y_i) , i \in [1, k]}$, of images $x_i$, and their associated labels $y_i$, and an unlabeled probe image $\hat{x}$,<br />
the one-shot learning task is to identify the true label of the probe image, $\hat{y}$, from the support set labels $ {y_i , i \in [1, k]} $:<br />
<br />
<br />
$\displaystyle \hat{y} = arg \max_{y}$ $P(y | \hat{x}, S)$<br />
<br />
<br />
We assume that the image labels $y_i$ are represented using a one-hot encoding and that $P(y|\hat{x}, S)$ is parameterised by a DNN, allowing us to leverage the ability of deep networks to learn powerful representations.<br />
<br />
== Inception Networks ==<br />
<br />
A probe image $\hat{x}$ is given the label of the nearest neighbour from the<br />
support set:<br />
<br />
$\hat{y} = y$<br />
<br />
$(x, y) = \displaystyle arg \min_{(x_i,y_i) \in S} d(h(x_i), h(\hat{x})) $<br />
<br />
where d is a distance function.<br />
<br />
The function h is parameterized by Inception – one of the best performing ImageNet classification models. Specifically, h returns features from the last layer (the softmax input) of a pre-trained Inception classifier. With these features as input and cosine distance as the distance function, the classifier in achieves 87.6% accuracy on one-shot classification on the ImageNet dataset (Vinyals et al., 2016). We call the Inception classifier together with the nearest-neighbor component the Inception Baseline (IB) model.<br />
<br />
<br />
== Matching Networks ==<br />
<br />
MNs (Vinyals et al.,2016) are neural network architectures with state-of-the-art one shot learning performance on ImageNet (93.2% one-shot labelling accuracy).<br />
MNs are trained to assign label $\hat{y}$ to probe image $\hat{x}$ using an attention mechanism $a$ acting on image embeddings stored in the support set S:<br />
\begin{align*}<br />
a(\hat{x},x_i)=\frac{e^{d(f(\hat{x},S), g(x_i,S))}}{\sum_{j}e^{d(f(\hat{x},S), g(x_j,S))}},<br />
\end{align*}<br />
<br />
<br />
where d is a cosine distance and where f and g provide context-dependent embeddings of $\hat{x}$ and $x_i$ (with contextS). The embedding $g(x_i, S)$ is a bi-directional LSTM (Hochreiter & Schmidhuber, 1997) with the support set S provided as an input sequence. The embedding $f(\hat{x}, S)$ is an LSTM with a read-attention mechanism operating over the entire embedded support set. The input to the LSTM is given by the penultimate layer features of a pre-trained deep convolutional network, specifically Inception.<br />
<br />
To train MNs we proceed as follows:<br />
<br />
=== Training MN ===<br />
* Step 1: At each step of training, the model is given a small support set of images and associated labels. In addition to the support set, the model is fed an unlabeled probe image $\hat{x}$<br />
<br />
* Step 2: The model parameters are then updated to improve classification accuracy of the probe image $\hat{x}$ given the support set. Parameters are updated using stochastic gradient descent with a learning rate of 0.1<br />
<br />
* Step 3: After each update, the labels ${(y_i, i \in [1, k]}$ in the training set are randomly re-assigned to new image classes (the label indices are randomly permuted,<br />
but the image labels are not changed). This is a critical step. It prevents MNs from learning a consistent mapping between a category and a label. Usually, in classification, this is what we want, but in one-shot learning we want to train our model for classification after viewing a single in-class example from the support set. <br />
<br />
The objective function used is:<br />
<br />
\begin{align*}<br />
L=E_{C\sim T}\biggr[E_{S \sim C, B \sim C}\bigr[\sum_{(x,y)\in B}\log P(y|x,S)\bigr]\biggr]<br />
\end{align*}<br />
<br />
<br />
where T is the set of all possible labelings of our classes, S is a support set sampled with a class labeling C ~ T and B is a batch of probe images and labels, also with the same randomly chosen class labeling as the support set.<br />
<br />
== Cognitive Biases ==<br />
A Cognitive bias, in simple terms, is any systematic deviation from logic. Wikipedia defines a Cognitive bias as "a systematic pattern of deviation from norm or rationality in judgment, whereby inferences about other people and situations may be drawn in an illogical fashion". Cognitive bias is a concept from developmental psychology which attempts to explain how children can extract meanings of words with very few examples, similar to the concept of one-shot learning discussed above. The theory, as explained by the authors, is that humans form biases that allow them to eliminate many potential hypotheses about word meaning where the amount of data available is insufficient for this purpose. These include:<br />
* Whole object bias<br />
* Taxonomic bias<br />
* Mutual exclusivity bias<br />
* Shape bias<br />
A more complete list of cognitive biases is given by [[#References|(Bloom, 2000)]]. The bias the authors investigate in this paper is the shape bias, which denotes a tendency to assign the same name to similarly shaped items rather than to items with similar color, texture, or size.<br />
<br />
= Methodology =<br />
== Inductive Biases & Probe Data ==<br />
<br />
Inductive biases are those criteria which are artificially selected or learnt by the network as a classifying/distinguishing property.<br />
It has been observed that the biases that DNNs learnt are complex composite features. We, as researchers can take advantage of the fact that DNNs learnt complex distinguishing features by constructing probe data sets which particularly target on exposing a particular bias that a DNN might have. <br />
<br />
* Step 1: Take a known composite feature which we suspect the DNNs are biased against. It is possible that the feature is not numerical but intuitive for human researchers to understand.<br />
* Step 2: Train the target model with an appropriate dataset.<br />
* Step 3: Transfer Learning: Use the pre-trained model with a new data set which is curated to contain data to prove/disprove the existence of the bias<br />
* Step 4: Model/Decide on a function which quantifies the bias under study<br />
* Step 5: Measure the bias with the bias function<br />
<br />
Then it is possible to deduce whether the DNN uses the feature to solve the task by the value of bias function.<br />
<br />
== Data Sets Used ==<br />
<br />
* Training Set: ImageNet<br />
* Test Set:<br />
** The Cognitive Psychology Probe Data (CogPsyc data) that is used consists of 150 images of objects. The images are arranged in triples consisting of a probe image, a shape-match image (that matches the probe in colour but not shape), and a color-match image (that matches the probe in shape but not colour). In the dataset there are 10 triples, each shown on 5 different backgrounds, giving a total of 50 triples. [[File:CogPsy.PNG|center|350px]]<br />
** A real-world dataset consisting of 90 images of objects (30 triples) collected using Google Image Search. The images are arranged in triples consisting of a probe, a shape-match and a colour-match.<br />
<br />
= Experiments =<br />
== Evaluation Criteria ==<br />
<br />
* For a given probe image $\hat{x}$, we loaded the shape-match image $x_s$ and corresponding label $y_s$, along with the colour-match image $x_c$ and corresponding label $y_c$ into memory, as the support set $S = \{(x_s, ys), (x_c, y_c)\}$<br />
* Calculate $\hat{y}$<br />
* The model assigns either $y_c$ or $y_s$ to the probe image.<br />
* To estimate the shape bias Bs, calculate the proportion of shape labels assigned to the probe: $B_s = E(\delta(\hat{y} - y_s))$ <br />
where E is an expectation across probe images and $\delta$ is the Dirac delta function.<br />
<br />
== Experiment 1: Shape bias statistics in Inception Baseline: ==<br />
* Shape bias of IB to be $B_s = 0.68$. Similarly, the shape bias of IB using our real-world dataset was $B_s = 0.97$. Together, these results strongly suggest that IB trained on ImageNet has a stronger bias towards shape than colour<br />
<br />
== Experiment 2: Shape bias statistics in Matching Network: ==<br />
* They found that MNs have a shape of bias $B_s = 0.7$ using the CogPsyc dataset and a bias of $Bs = 1$ using the real-world dataset. Once again, these results suggest that MNs trained seeding from Inception using ImageNet has a stronger bias towards shape than colour.<br />
<br />
== Experiment 3: Shape bias statistics between and across models: ==<br />
<br />
The authors extended the shape bias analysis to calculate the shape bias in a population of IB models and in a population of MN models with different random initialization<br />
<br />
=== Dependence on the initialization of parameters: ===<br />
<br />
[[File:3.1.PNG|right|250px]] A strong variability was observed when variation in the initial values of the parameters. For the CogPsyc dataset, the average shape bias was $B_s = 0.628$ with standard deviation $\sigma B_s = 0.049$ at the end of training and for the real-world dataset the average shape bias was $B_s = 0:958$ with $\sigma B_s = 0.037$.<br />
<br />
=== Dependence of shape bias on model performance: ===<br />
<br />
For the CogPsych dataset, the correlation between bias and classification accuracy was $\rho = 0.15$, and for the real world dataset, correlation between bias and classification accuracy was $\rho = -0.06$. This would be evident since the accuracy of the models remained nearly constant when the initialization parameters varied whereas the shape bias tended to vary a lot, hence highlighting the lack of correlation amongst them. <br />
<br />
=== Emergence of shape bias during training: ===<br />
The shape bias spiked to a large value very early. <br />
<br />
=== Variation of shape bias within models & across models: ===<br />
With different initialization parameters, the shape bias varied a lot within IB during training while the shape bias did not fluctuate during the training of MN. It was found that the MN inherits the shape bias of the IB which seeded its embeddings and thereafter, the shape bias remained constant throughout training. It is important to note that the output of the penultimate layer of the Inception was not fine tuned when it was pipelined to the MN. This was to ensure that the MN properties were independent of the IB model properties. [[File:3.3.PNG|center|250px]] [[File:3.4.PNG|center|250px]]<br />
<br />
= Learnings, Inferences & Implications =<br />
* Both the Inception Baseline and the Matching Network exhibit strong shape bias when trained on ImageNet. Researchers who use Inception & MN DNNs can now use this fact as a consideration for their application while using pre-trained models for new datasets. If it is known before hand that the new data set is strongly classifiable through a color bias, then they would either want to defer using the pre-trained models or explore methods to decrease/remove the strong shape bias.<br />
<br />
* There exists a high variability in the shape bias with the variation in the initialization parameters. This is an important finding since it uncovers the fact that the same architecture which exhibit similar accuracy in predictions can display a variety of shape bias just with different initialization parameters. Researchers can explore methods of tuning the random initialization such that the models start out with a low shape bias without compromising the accuracy of the model.<br />
<br />
* MNs inherit the shape bias which is seeded to it by the Inception Network's input embedding. This is also another fact which researchers & practitioners should be careful about. When using cascaded or pipelined heterogeneous architectures, the models downstream tend to inherit/become/are fed with the properties/biases of the models upstream. This may be desirable or undesirable according to the application, but it is important to be aware of its presence.<br />
<br />
* The biases under consideration are the property of the collection of the architecture, the dataset and the optimization procedure. Hence in order to increase or decrease the effect of a particular bias, one or more of the mentioned factors must be adjusted/tuned/changed.<br />
<br />
* The fact that a high shape bias emerged in the early epochs with less variability in further epochs can be thought of analogous to the biases that humans develop at an infancy which gets fortified as they age.<br />
<br />
=== Modeling human word learning === <br />
The authors note how there have been several previous attempts to model human word learning in the field of cognitive science. A major shortcoming of these works is that none of the models are capable of one-shot word learning when presented with real-world images. Recognizing this and given the success of MNs in this work, the authors propose MNs as a computational-level account of one-shot word learning. <br />
In the paper, it is discussed how shape bias increases dramatically during early training of the model. This agrees with what is observed in many psychological studies: older children show bias more than younger children. <br />
<br />
= Conclusion, Future Work and Open questions =<br />
<br />
* Just as cognitive psychology exposes the shape bias observed in this experiment, we should try to uncover other biases as well using multiple approaches<br />
* Study the underlying mechanisms which cause biases such as shape bias in DNNs<br />
* Research into various methods of probing and creating probe data sets which can be used to test architectures for various biases<br />
* Exploration into a research field called Artificial Cognitive Psychology which focuses on probing how DNN architectures can be understood further using known behaviors of the human brain<br />
<br />
= References =<br />
<br />
* Ritter, Samuel & G. T. Barrett, David & Santoro, Adam & M. Botvinick, Matt. (2017). Cognitive Psychology for Deep Neural Networks: A Shape Bias Case Study<br />
<br />
* Vinyals, Oriol, Blundell, Charles, Lillicrap, Timothy, Kavukcuoglu, Koray, and Wierstra, Daan. Matching networks for one shot learning. arXiv preprint arXiv:1606.04080, 2016.<br />
<br />
* Bloom, P. (2000). How children learn the meanings of words. The MIT Press.<br />
<br />
* https://www.slideshare.net/KazukiFujikawa/matching-networks-for-one-shot-learning-71257100<br />
<br />
* https://deepmind.com/blog/cognitive-psychology/<br />
<br />
* https://hacktilldawn.com/2016/09/25/inception-modules-explained-and-implemented/</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/Cognitive_Psychology_For_Deep_Neural_Networks:_A_Shape_Bias_Case_Study&diff=29761STAT946F17/Cognitive Psychology For Deep Neural Networks: A Shape Bias Case Study2017-11-09T00:43:36Z<p>Nkadiyar: /* Matching Networks */</p>
<hr />
<div>= Introduction =<br />
<br />
The recent burgeon on the use of Deep Neural Networks (DNNs) have resulted in giant leaps of accuracy in prediction. They are also being used to solve a variety of complex tasks which earlier methodologies have struggled to excel in.<br />
<br />
While it is all good to see incredibly high accuracy as a result of the use of DNN, we must begin to question why they perform so well. It has become an interesting field of study to actually represent the features/feature maps or interpret the meaning of the learnt values in a DNN's hidden layers. Currently we treat models of DNNs as black boxes which we practically tune the tweakable parameters like number of layers, number of units in each layer, number & size of feature maps(in case of CNN) etc. The opacity created by the lack of an intuitive representation of the internal learnt parameters of DNNs hinders both basic research as well as its application to real world problems.<br />
<br />
Recent pushes have aimed to better understand DNNs: tailor-made loss functions and architectures produce more interpretable features (Higgins et al., 2016; Raposo et al., 2017) while output-behavior analyses unveil previously opaque operations of these networks (Karpathy et al., 2015). Parallel to this work, neuroscience-inspired methods such as activation visualization (Li et al., 2015), ablation analysis (Zeiler & Fergus, 2014) and activation maximization (Yosinski et al., 2015) have also been applied<br />
<br />
This paper aims to provide another methodology to attempt to decipher & better understand how DNNs solve a particular task. This methodology was inspired by psychological concepts to test whether the DNN's were able to make accurate predictions with biases similar to that the human mind makes.<br />
<br />
Research in developmental psychology shows that when learning new words, humans tend to assign the same name to similarly shaped items rather than to items with similar color, texture, or size. This bias/knowledge tend to be forged into the brains of humans and humans then take this forward to easily associate these shapes with new objects they have not seen before.<br />
<br />
The authors of this paper try to test if DNNs behave similarly in one-shot learning applications. They attempt to prove that when the models of state-of-the-art DNNs are used to learn objects from images, they exhibit a stronger shape bias than a color bias. To emulate the human brain, they use the parameters of pre-trained DNN models and use this to perform one-shot learning on a new data set with different labels.<br />
<br />
= Background =<br />
== One Shot Learning ==<br />
One-shot learning is an object categorization problem in computer vision. Whereas most machine learning based object categorization algorithms require training on hundreds or thousands of images and very large datasets, one-shot learning aims to learn information about object categories from one, or only a few, training images.<br />
<br />
The one-shot word learning task is to label a novel data example $\hat{x}$ (e.g. a novel probe image) with a novel class label $\hat{y}$ (e.g. a new word) after only a single example.<br />
<br />
More specifically, given a support set $S = {(x_i, y_i) , i \in [1, k]}$, of images $x_i$, and their associated labels $y_i$, and an unlabeled probe image $\hat{x}$,<br />
the one-shot learning task is to identify the true label of the probe image, $\hat{y}$, from the support set labels $ {y_i , i \in [1, k]} $:<br />
<br />
<br />
$\displaystyle \hat{y} = arg \max_{y}$ $P(y | \hat{x}, S)$<br />
<br />
<br />
We assume that the image labels $y_i$ are represented using a one-hot encoding and that $P(y|\hat{x}, S)$ is parameterised by a DNN, allowing us to leverage the ability of deep networks to learn powerful representations.<br />
<br />
== Inception Networks ==<br />
<br />
A probe image $\hat{x}$ is given the label of the nearest neighbour from the<br />
support set:<br />
<br />
$\hat{y} = y$<br />
<br />
$(x, y) = \displaystyle arg \min_{(x_i,y_i) \in S} d(h(x_i), h(\hat{x})) $<br />
<br />
where d is a distance function.<br />
<br />
The function h is parameterized by Inception – one of the best performing ImageNet classification models. Specifically, h returns features from the last layer (the softmax input) of a pre-trained Inception classifier. With these features as input and cosine distance as the distance function, the classifier in achieves 87.6% accuracy on one-shot classification on the ImageNet dataset (Vinyals et al., 2016). We call the Inception classifier together with the nearest-neighbor component the Inception Baseline (IB) model.<br />
<br />
<br />
== Matching Networks ==<br />
<br />
MNs (Vinyals et al.,2016) are neural network architectures with state-of-the-art one shot learning performance on ImageNet (93.2% one-shot labelling accuracy).<br />
MNs are trained to assign label $\hat{y}$ to probe image $\hat{x}$ using an attention mechanism $a$ acting on image embeddings stored in the support set S:<br />
\begin{align*}<br />
a(\hat{x},x_i)=\frac{e^{d(f(\hat{x},S), g(x_i,S))}}{\sum_{j}e^{d(f(\hat{x},S), g(x_j,S))}},<br />
\end{align*}<br />
<br />
<br />
where d is a cosine distance and where f and g provide context-dependent embeddings of $\hat{x}$ and $x_i$ (with contextS). The embedding $g(x_i, S)$ is a bi-directional LSTM (Hochreiter & Schmidhuber, 1997) with the support set S provided as an input sequence. The embedding $f(\hat{x}, S)$ is an LSTM with a read-attention mechanism operating over the entire embedded support set. The input to the LSTM is given by the penultimate layer features of a pre-trained deep convolutional network, specifically Inception.<br />
<br />
To train MNs we proceed as follows:<br />
<br />
=== Training MN ===<br />
* Step 1: At each step of training, the model is given a small support set of images and associated labels. In addition to the support set, the model is fed an unlabeled probe image $\hat{x}$<br />
<br />
* Step 2: The model parameters are then updated to improve classification accuracy of the probe image $\hat{x}$ given the support set. Parameters are updated using stochastic gradient descent with a learning rate of 0.1<br />
<br />
* Step 3: After each update, the labels ${(y_i, i \in [1, k]}$ in the training set are randomly re-assigned to new image classes (the label indices are randomly permuted,<br />
but the image labels are not changed). This is a critical step. It prevents MNs from learning a consistent mapping between a category and a label. Usually, in classification, this is what we want, but in one-shot learning we want to train our model for classification after viewing a single in-class example from the support set. <br />
<br />
The objective function used is:<br />
<br />
\begin{align*}<br />
L=E_{C\sim T}\biggr[E_{S \sim C, B \sim C}\bigr[\sum_{(x,y)\in B}\log P(y|x,S)\bigr]\biggr]<br />
\end{align*}<br />
<br />
<br />
where T is the set of all possible labelings of our classes, S is a support set sampled with a class labeling C ~ T and B is a batch of probe images and labels, also with the same randomly chosen class labeling as the support set.<br />
<br />
== Cognitive Biases ==<br />
Cognitive bias is a concept from developmental psychology which attempts to explain how children can extract meanings of words with very few examples, similar to the concept of one-shot learning discussed above. The theory, as explained by the authors, is that humans form biases that allow them to eliminate many potential hypotheses about word meaning where the amount of data available is insufficient for this purpose. These include:<br />
* Whole object bias<br />
* Taxonomic bias<br />
* Mutual exclusivity bias<br />
* Shape bias<br />
A more complete list of cognitive biases is given by [[#References|(Bloom, 2000)]]. The bias the authors investigate in this paper is the shape bias, which denotes a tendency to assign the same name to similarly shaped items rather than to items with similar color, texture, or size.<br />
<br />
= Methodology =<br />
== Inductive Biases & Probe Data ==<br />
<br />
Inductive biases are those criteria which are artificially selected or learnt by the network as a classifying/distinguishing property.<br />
It has been observed that the biases that DNNs learnt are complex composite features. We, as researchers can take advantage of the fact that DNNs learnt complex distinguishing features by constructing probe data sets which particularly target on exposing a particular bias that a DNN might have. <br />
<br />
* Step 1: Take a known composite feature which we suspect the DNNs are biased against. It is possible that the feature is not numerical but intuitive for human researchers to understand.<br />
* Step 2: Train the target model with an appropriate dataset.<br />
* Step 3: Transfer Learning: Use the pre-trained model with a new data set which is curated to contain data to prove/disprove the existence of the bias<br />
* Step 4: Model/Decide on a function which quantifies the bias under study<br />
* Step 5: Measure the bias with the bias function<br />
<br />
Then it is possible to deduce whether the DNN uses the feature to solve the task by the value of bias function.<br />
<br />
== Data Sets Used ==<br />
<br />
* Training Set: ImageNet<br />
* Test Set:<br />
** The Cognitive Psychology Probe Data (CogPsyc data) that is used consists of 150 images of objects. The images are arranged in triples consisting of a probe image, a shape-match image (that matches the probe in colour but not shape), and a color-match image (that matches the probe in shape but not colour). In the dataset there are 10 triples, each shown on 5 different backgrounds, giving a total of 50 triples. [[File:CogPsy.PNG|center|350px]]<br />
** A real-world dataset consisting of 90 images of objects (30 triples) collected using Google Image Search. The images are arranged in triples consisting of a probe, a shape-match and a colour-match.<br />
<br />
= Experiments =<br />
== Evaluation Criteria ==<br />
<br />
* For a given probe image $\hat{x}$, we loaded the shape-match image $x_s$ and corresponding label $y_s$, along with the colour-match image $x_c$ and corresponding label $y_c$ into memory, as the support set $S = \{(x_s, ys), (x_c, y_c)\}$<br />
* Calculate $\hat{y}$<br />
* The model assigns either $y_c$ or $y_s$ to the probe image.<br />
* To estimate the shape bias Bs, calculate the proportion of shape labels assigned to the probe: $B_s = E(\delta(\hat{y} - y_s))$ <br />
where E is an expectation across probe images and $\delta$ is the Dirac delta function.<br />
<br />
== Experiment 1: Shape bias statistics in Inception Baseline: ==<br />
* Shape bias of IB to be $B_s = 0.68$. Similarly, the shape bias of IB using our real-world dataset was $B_s = 0.97$. Together, these results strongly suggest that IB trained on ImageNet has a stronger bias towards shape than colour<br />
<br />
== Experiment 2: Shape bias statistics in Matching Network: ==<br />
* They found that MNs have a shape of bias $B_s = 0.7$ using the CogPsyc dataset and a bias of $Bs = 1$ using the real-world dataset. Once again, these results suggest that MNs trained seeding from Inception using ImageNet has a stronger bias towards shape than colour.<br />
<br />
== Experiment 3: Shape bias statistics between and across models: ==<br />
<br />
The authors extended the shape bias analysis to calculate the shape bias in a population of IB models and in a population of MN models with different random initialization<br />
<br />
=== Dependence on the initialization of parameters: ===<br />
<br />
[[File:3.1.PNG|right|250px]] A strong variability was observed when variation in the initial values of the parameters. For the CogPsyc dataset, the average shape bias was $B_s = 0.628$ with standard deviation $\sigma B_s = 0.049$ at the end of training and for the real-world dataset the average shape bias was $B_s = 0:958$ with $\sigma B_s = 0.037$.<br />
<br />
=== Dependence of shape bias on model performance: ===<br />
<br />
For the CogPsych dataset, the correlation between bias and classification accuracy was $\rho = 0.15$, and for the real world dataset, correlation between bias and classification accuracy was $\rho = -0.06$. This would be evident since the accuracy of the models remained nearly constant when the initialization parameters varied whereas the shape bias tended to vary a lot, hence highlighting the lack of correlation amongst them. <br />
<br />
=== Emergence of shape bias during training: ===<br />
The shape bias spiked to a large value very early. <br />
<br />
=== Variation of shape bias within models & across models: ===<br />
With different initialization parameters, the shape bias varied a lot within IB during training while the shape bias did not fluctuate during the training of MN. It was found that the MN inherits the shape bias of the IB which seeded its embeddings and thereafter, the shape bias remained constant throughout training. It is important to note that the output of the penultimate layer of the Inception was not fine tuned when it was pipelined to the MN. This was to ensure that the MN properties were independent of the IB model properties. [[File:3.3.PNG|center|250px]] [[File:3.4.PNG|center|250px]]<br />
<br />
= Learnings, Inferences & Implications =<br />
* Both the Inception Baseline and the Matching Network exhibit strong shape bias when trained on ImageNet. Researchers who use Inception & MN DNNs can now use this fact as a consideration for their application while using pre-trained models for new datasets. If it is known before hand that the new data set is strongly classifiable through a color bias, then they would either want to defer using the pre-trained models or explore methods to decrease/remove the strong shape bias.<br />
<br />
* There exists a high variability in the shape bias with the variation in the initialization parameters. This is an important finding since it uncovers the fact that the same architecture which exhibit similar accuracy in predictions can display a variety of shape bias just with different initialization parameters. Researchers can explore methods of tuning the random initialization such that the models start out with a low shape bias without compromising the accuracy of the model.<br />
<br />
* MNs inherit the shape bias which is seeded to it by the Inception Network's input embedding. This is also another fact which researchers & practitioners should be careful about. When using cascaded or pipelined heterogeneous architectures, the models downstream tend to inherit/become/are fed with the properties/biases of the models upstream. This may be desirable or undesirable according to the application, but it is important to be aware of its presence.<br />
<br />
* The biases under consideration are the property of the collection of the architecture, the dataset and the optimization procedure. Hence in order to increase or decrease the effect of a particular bias, one or more of the mentioned factors must be adjusted/tuned/changed.<br />
<br />
* The fact that a high shape bias emerged in the early epochs with less variability in further epochs can be thought of analogous to the biases that humans develop at an infancy which gets fortified as they age.<br />
<br />
=== Modeling human word learning === <br />
The authors note how there have been several previous attempts to model human word learning in the field of cognitive science. A major shortcoming of these works is that none of the models are capable of one-shot word learning when presented with real-world images. Recognizing this and given the success of MNs in this work, the authors propose MNs as a computational-level account of one-shot word learning. <br />
In the paper, it is discussed how shape bias increases dramatically during early training of the model. This agrees with what is observed in many psychological studies: older children show bias more than younger children. <br />
<br />
= Conclusion, Future Work and Open questions =<br />
<br />
* Just as cognitive psychology exposes the shape bias observed in this experiment, we should try to uncover other biases as well using multiple approaches<br />
* Study the underlying mechanisms which cause biases such as shape bias in DNNs<br />
* Research into various methods of probing and creating probe data sets which can be used to test architectures for various biases<br />
* Exploration into a research field called Artificial Cognitive Psychology which focuses on probing how DNN architectures can be understood further using known behaviors of the human brain<br />
<br />
= References =<br />
<br />
* Ritter, Samuel & G. T. Barrett, David & Santoro, Adam & M. Botvinick, Matt. (2017). Cognitive Psychology for Deep Neural Networks: A Shape Bias Case Study<br />
<br />
* Vinyals, Oriol, Blundell, Charles, Lillicrap, Timothy, Kavukcuoglu, Koray, and Wierstra, Daan. Matching networks for one shot learning. arXiv preprint arXiv:1606.04080, 2016.<br />
<br />
* Bloom, P. (2000). How children learn the meanings of words. The MIT Press.<br />
<br />
* https://www.slideshare.net/KazukiFujikawa/matching-networks-for-one-shot-learning-71257100<br />
<br />
* https://deepmind.com/blog/cognitive-psychology/<br />
<br />
* https://hacktilldawn.com/2016/09/25/inception-modules-explained-and-implemented/</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Dialog-based_Language_Learning&diff=29760Dialog-based Language Learning2017-11-09T00:41:06Z<p>Nkadiyar: /* Training strategies */</p>
<hr />
<div>This page is a summary for NIPS 2016 paper <i>Dialog-based Language Learning</i> [1].<br />
==Introduction==<br />
One of the ways humans learn language, especially second language or language learning by students, is by communication and getting its feedback. However, most existing research in Natural Language Understanding has focused on supervised learning from fixed training sets of labeled data. This kind of supervision is not realistic of how humans learn, where language is both learned by, and used for, communication. When humans act in dialogs (i.e., make speech utterances) the feedback from other human’s responses contain very rich information. This is perhaps most pronounced in a student/teacher scenario where the teacher provides positive feedback for successful communication and corrections for unsuccessful ones. <br />
<br />
This paper is about dialog-based language learning, where supervision is given naturally and implicitly in the response of the dialog partner during the conversation. This paper is a step towards the ultimate goal of being able to develop an intelligent dialog agent that can learn while conducting conversations. Specifically this paper explores whether we can train machine learning models to learn from dialog.<br />
<br />
===Contributions of this paper===<br />
*Introduce a set of tasks that model natural feedback from a teacher and hence assess the feasibility of dialog-based language learning. <br />
*Evaluated some baseline models on this data and compared them to standard supervised learning. <br />
*Introduced a novel forward prediction model, whereby the learner tries to predict the teacher’s replies to its actions, which yields promising results, even with no reward signal at all<br />
<br />
==Background on Memory Networks==<br />
<br/><br />
[[File:DB F1.png|center|700px]]<br />
<br/><br />
A memory network combines learning strategies from the machine learning literature with a memory component that can be read and written to.<br />
<br />
The high-level view of a memory network is as follows:<br />
*There is a memory, $m$, an indexed array of objects (e.g. vectors or arrays of strings).<br />
*An input feature map $I$, which converts the incoming input to the internal feature representation<br />
*A generalization component $G$ which updates old memories given the new input. <br />
*An output feature map $O$, which produces a new output in the feature representation space given the new input and the current memory state.<br />
*A response component $R$ which converts the output into the response format desired – for example, a textual response or an action.<br />
<br />
$I$, $G$, $O$ and $R$ can all potentially be learned components and make use of any ideas from the existing machine learning literature.<br />
<br />
In question answering systems for example, the components may be instantiated as follows:<br />
*$I$ can make use of standard pre-processing such as parsing, coreference, and entity resolution. It could also encode the input into an internal feature representation by converting from text to a sparse or dense feature vector.<br />
*The simplest form of $G$ is to introduce a function $H$ which maps the internal feature representation produced by I to an individual memory slot, and just updates the memory at $H(I(x))$.<br />
*$O$ Reads from memory and performs inference to deduce the set of relevant memories needed to perform a good response.<br />
*$R$ would produce the actual wording of the question answer based on the memories found by $O$. For example, $R$ could be an RNN conditioned on the output of $O$<br />
<br />
When the components $I$,$G$,$O$, & $R$ are neural networks, the authors describe the resulting system as a <b>Memory Neural Network (MemNN)</b>. They build a MemNN for QA (question answering) problems and compare it to RNNs (Recurrent Neural Network) and LSTMs (Long Short Term Memory RNNs) and find that it gives superior performance.<br />
<br />
[[File:DB_F2.png|center|800px]]<br />
<br />
==Related Work==<br />
<br />
'''Usefulness of feedback in language learning:''' Social interaction and natural infant directed conversations are shown to useful for language learning[2]. Several studies[3][4][5][6] have shown that feedback is especially useful in second language learning and learning by students.<br />
<br />
'''Supervised learning from dialogs using neural models:''' Neural networks has been used for response generation that can be trained end to end on large quantities of unstructured Twitter conversations[7]. However this does not incorporate feedback from dialog partner during real time conversation<br />
<br />
'''Reinforcement learning:''' Reinforcement learning works on dialogs[8][9], often consider reward as the feedback model rather than exploiting the dialog feedback per se.<br />
<br />
'''Forward prediction models:''' Although forward prediction models, have been used in other applications like learning eye-tracking[10], controlling robot arms[11] and vehicles[12], it has not been used for dialog.<br />
<br />
==Dialog-based Supervision tasks==<br />
For testing their models, the authors chose two datasets (i) the single supporting fact problem from the bAbI datasets [13] which consists of short stories from a simulated world followed by questions; and (ii) the MovieQA dataset [14] which is a large-scale dataset (∼ 100k questions over ∼ 75k entities) based on questions with answers in the open movie database (OMDb)<br />
<br />
However, since these datasets were not designed to model the supervision from dialogs, the authors modified them to create 10 supervision task types on these datasets(Fig 3).<br />
<br />
[[File:DB F3.png|center|700px]]<br />
<br/><br />
<br />
*'''Task 1: Imitating an Expert Student''': The dialogs take place between a teacher and an expert student who gives semantically coherent answers. Hence, the task is for the learner to imitate that expert student, and become an expert themselves <br />
<br />
*'''Task 2: Positive and Negative Feedback:''' When the learner answers a question the teacher then replies with either positive or negative feedback. In the experiments the subsequent responses are variants of “No, that’s incorrect” or “Yes, that’s right”. In the datasets there are 6 templates for positive feedback and 6 templates for negative feedback, e.g. ”Sorry, that’s not it.”, ”Wrong”, etc. To distinguish the notion of positive from negative, an additional external reward signal that is not part of the text<br />
<br />
*'''Task 3: Answers Supplied by Teacher:''' The teacher gives positive and negative feedback as in Task 2, however when the learner’s answer is incorrect, the teacher also responds with the correction. For example if “where is Mary?” is answered with the incorrect answer “bedroom” the teacher responds “No, the answer is kitchen”’<br />
<br />
*'''Task 4: Hints Supplied by Teacher:''' The corrections provided by the teacher do not provide the exact answer as in Task 3, but only a useful hint. This setting is meant to mimic the real life occurrence of being provided only partial information about what you did wrong.<br />
<br />
*'''Task 5: Supporting Facts Supplied by Teacher:''' Another way of providing partial supervision for an incorrect answer is explored. Here, the teacher gives a reason (explanation) why the answer is wrong by referring to a known fact that supports the true answer that the incorrect answer may contradict. <br />
<br />
*'''Task 6: Partial Feedback:''' External rewards are only given some of (50% of) the time for correct answers, the setting is otherwise identical to Task 3. This attempts to mimic the realistic situation of some learning being more closely supervised (a teacher rewarding you for getting some answers right) whereas other dialogs have less supervision (no external rewards). The task attempts to assess the impact of such partial supervision.<br />
<br />
*'''Task 7: No Feedback:''' External rewards are not given at all, only text, but is otherwise identical to Tasks 3 and 6. This task explores whether it is actually possible to learn how to answer at all in such a setting.<br />
<br />
*'''Task 8: Imitation and Feedback Mixture:''' Combines Tasks 1 and 2. The goal is to see if a learner can learn successfully from both forms of supervision at once. This mimics a child both observing pairs of experts talking (Task 1) while also trying to talk (Task 2).<br />
<br />
*'''Task 9: Asking For Corrections:''' The learner will ask questions to the teacher about what it has done wrong. Task 9 tests one of the most simple instances, where asking “Can you help me?” when wrong obtains from the teacher the correct answer.<br />
<br />
*'''Task 10: Asking for Supporting Facts:''' A less direct form of supervision for the learner after asking for help is to receive a hint rather than the correct answer, such as “A relevant fact is John moved to the bathroom” when asking “Can you help me?”. This is thus related to the supervision in Task 5 except the learner must request help<br />
<br />
[[File:F4.png|center|700px]]<br />
<br/><br />
<br />
For each task a fixed policy is considered for performing actions (answering questions) which gets questions correct with probability $π_{acc}$ (i.e. the chance of getting the red text correct in Figs. 3 and 4). We thus can compare different learning algorithms for each task over different values of $π_{acc}$ (0.5, 0.1 and 0.01). In all cases a training, validation and test set is provided.<br />
<br />
==Learning models==<br />
This work evaluates four possible learning strategies for each of the 10 tasks: imitation learning, reward-based imitation, forward prediction, and a combination of reward-based imitation and forward prediction<br />
<br />
All of these approaches are evaluated with the same model architecture: an end-to-end memory network (MemN2N) [15], which has been used as a baseline model for exploring differing modes of learning.<br />
<br />
[[File:F5.png|center|700px]]<br />
<br/><br />
The input is the last utterance of the dialog, $x$, as well as a set of memories (context) (<math> c_1</math>, . . . , <math> c_n</math> ) which can encode both short-term memory, e.g. recent previous utterances and replies, and long-term memories, e.g. facts that could be useful for answering questions. The context inputs <math> c_i</math> are converted into vectors <math> m_i</math> via embeddings and are stored in the memory. The goal is to produce an output $\hat{a}$ by processing the input $x$ and using that to address and read from the memory, $m$, possibly multiple times, in order to form a coherent reply. In the figure the memory is read twice, which is termed multiple “hops” of attention. <br />
<br />
In the first step, the input $x$ is embedded using a matrix $A$ of size $d$ × $V$ where $d$ is the embedding dimension and $V$ is the size of the vocabulary, giving $q$ = $A$$x$, where the input $x$ is as a bag-of words vector. Each memory <math> c_i</math> is embedded using the same matrix, giving $m_i$ = $A$$c_i$ . The output of addressing and then reading from memory in the first hop is: <br />
<br />
[[File:eq1.png|center|400px]]<br />
<br />
Here, $p^{1}$ is a probability vector over the memories, and is a measure of how much the input and the memories match. The goal is to select memories relevant to the last utterance $x$, i.e. the most relevant have large values of $p^{1}_i$ . The output memory representation $o_1$ is then constructed using the weighted sum of memories, i.e. weighted by $p^{1}$ . The memory output is then added to the original input, <math> c_1</math> = <math> R_1</math>(<math> o_1</math> + $q$), to form the new state of the controller, where <math> R_1</math> is a $d$ × $d$ rotation matrix . The attention over the memory can then be repeated using <math> u_1</math> instead of $q$ as the addressing vector, yielding: <br />
<br />
[[File:eq2.png|center|400px]]<br />
<br />
The controller state is updated again with <math> u_2</math> = <math> R_2</math>(<math> o_2</math> + <math> u_1</math>), where <math> R_2</math> is another $d$ × $d$ matrix to be learnt. In a two-hop model the final output is then defined as: <br />
<br />
[[File:eq3.png|center|400px]]<br />
<br />
where there are $C$ candidate answers in $y$. In our experiments $C$ is the set of actions that occur in the training set for the bAbI tasks, and for MovieQA it is the set of words retrieved from the KB.<br />
<br />
==Training strategies==<br />
*'''1. Imitation Learning'''<br />
This approach involves simply imitating one of the speakers in observed dialogs. Examples arrive as $(x, c, a)$ triples, where $a$ is (assumed to be) a good response to the last utterance $x$ given context $c$. Here, the whole memory network model defined above is trained using stochastic gradient descent by minimizing a standard cross-entropy loss between $\hat{a}$ and the label $a$<br />
<br />
*'''2. Reward-based Imitation''' <br />
If some actions are poor choices, then one does not want to repeat them, that is we shouldn’t treat them as a supervised objective. Here, positive reward is only obtained immediately after (some of) the correct actions, or else is zero. Only apply imitation learning on the rewarded actions. The rest of the actions are simply discarded from the training set. <br />
<br />
*'''3. Forward Prediction''' <br />
The aim is, given an utterance $x$ from speaker 1 and an answer a by speaker 2 (i.e., the learner), to predict $x^{¯}$, the response to the answer from speaker 1. That is, in general to predict the changed state of the world after action $a$, which in this case involves the new utterance $x^{¯}$.<br />
<br />
[[File:F6.png|center|700px]]<br />
<br/><br />
As shown in Figure (b), this is achieved by chopping off the final output from the original network of Fig (a) and replace it with some additional layers that compute the forward prediction. The first part of the network remains exactly the same and only has access to input x and context c, just as before. The computation up to $u_2$ = $R_2$($o_2$ + $u_1$) is thus exactly the same as before. <br />
<br />
Then perform another “hop” of attention but over the candidate answers rather than the memories. The information of which action (candidate) was actually selected in the dialog (i.e. which one is a) is also incorporated which is crucial. After this “hop”, the resulting state of the controller is then used to do the forward prediction.<br />
<br />
Concretely, we compute: <br />
<br />
[[File:eq4.png|center|550px]]<br />
<br />
where $β^{*}$ is a d-dimensional vector, that is also learnt, that represents in the output $o_3$ the action that was actually selected. After obtaining $o_3$, the forward prediction is then computed as: <br />
<br />
[[File:eq5.png|center|500px]]<br />
<br />
where $u_3$ = $R_3$($o_3$ + $u_2$). That is, it computes the scores of the possible responses to the answer a over $C^{¯}$ possible candidates.<br />
<br />
<br />
*'''4. Reward-based Imitation + Forward Prediction'''<br />
As the reward-based imitation learning uses the architecture of Fig (a), and forward prediction uses the same architecture but with the additional layers of Fig (b), we can learn jointly with both strategies. This is a powerful combination as it makes use of reward signal when available and the dialog feedback when the reward signal is not available.<br />
<br />
==Experiments==<br />
<br />
Experiments were conducted on the two test datasets - babi and MovieQA. For each task, a fixed policy is considered for performing actions (answering questions) which gets questions correct with probability $π_{acc}$. This helps to compare the different training strategies described earlier over each task for different values of $π_{acc}$. Hyperparameters for all methods are optimized on the validation sets.<br />
<br/><br/><br />
[[File:DB F7.png|center|800px]]<br />
<br/><br />
The following results are observed by the authors:<br />
*Imitation learning, ignoring rewards, is a poor learning strategy when imitating inaccurate answers, e.g. for $π_{acc}$ < 0.5. For imitating an expert however (Task 1) it is hard to beat. <br />
*Reward-based imitation (RBI) performs better when rewards are available, particularly in Table 1, but also degrades when they are too sparse e.g. for πacc = 0.01.<br />
*Forward prediction (FP) is more robust and has stable performance at different levels of πacc. However as it only predicts answers implicitly and does not make use of rewards it is outperformed by RBI on several tasks, notably Tasks 1 and 8 (because it cannot do supervised learning) and Task 2 (because it does not take advantage of positive rewards).<br />
*FP makes use of dialog feedback in Tasks 3-5 whereas RBI does not. This explains why FP does better with useful feedback (Tasks 3-5) than without (Task 2), whereas RBI cannot.<br />
*Supplying full answers (Task 3) is more useful than hints (Task 4) but hints still help FP more than just yes/no answers without extra information (Task 2).<br />
*When positive feedback is sometimes missing (Task 6) RBI suffers especially in Table 1. FP does not as it does not use this feedback.<br />
*One of the most surprising results of our experiments is that FP performs well overall, given that it does not use feedback, which we will attempt to explain subsequently. This is particularly evident on Task 7 (no feedback) where RBI has no hope of succeeding as it has no positive examples. FP on the other hand learns adequately.<br />
*Tasks 9 and 10 are harder for FP as the question is not immediately before the feedback.<br />
*Combining RBI and FP ameliorates the failings of each, yielding the best overall results<br />
<br />
==Future work==<br />
* Any reply in a dialog can be seen as feedback, and should be useful for learning. Evaluate if forward prediction, and the other approaches in this paper, work there too. <br />
* Develop further evaluation methodologies to test how the models presented here work in more complex settings where actions that are made lead to long-term changes in the environment and delayed rewards, i.e. extending to the reinforcement learning setting, and to full language generation. <br />
* How dialog-based feedback could also be used as a medium to learn non-dialog based skills, e.g. natural language dialog for completing visual or physical tasks.<br />
<br />
==References==<br />
# Jason Weston. Dialog-based Language Learning. NIPS, 2016.<br />
# P. K. Kuhl. Early language acquisition: cracking the speech code. Nature reviews neuroscience, 5(11): 831–843, 2004.<br />
# M. A. Bassiri. Interactional feedback and the impact of attitude and motivation on noticing l2 form. English Language and Literature Studies, 1(2):61, 2011.<br />
# R. Higgins, P. Hartley, and A. Skelton. The conscientious consumer: Reconsidering the role of assessment feedback in student learning. Studies in higher education, 27(1):53–64, 2002.<br />
# A. S. Latham. Learning through feedback. Educational Leadership, 54(8):86–87, 1997.<br />
# M. G. Werts, M. Wolery, A. Holcombe, and D. L. Gast. Instructive feedback: Review of parameters and effects. Journal of Behavioral Education, 5(1):55–75, 1995.<br />
# A. Sordoni, M. Galley, M. Auli, C. Brockett, Y. Ji, M. Mitchell, J.-Y. Nie, J. Gao, and B. Dolan. A neural network approach to context-sensitive generation of conversational responses. Proceedings of NAACL, 2015.<br />
# V. Rieser and O. Lemon. Reinforcement learning for adaptive dialogue systems: a data-driven methodology for dialogue management and natural language generation. Springer Science & Business Media, 2011.<br />
# J. Schatzmann, K. Weilhammer, M. Stuttle, and S. Young. A survey of statistical user simulation techniques for reinforcement-learning of dialogue management strategies. The knowledge engineering review, 21(02):97–126, 2006.<br />
# J. Schmidhuber and R. Huber. Learning to generate artificial fovea trajectories for target detection. International Journal of Neural Systems, 2(01n02):125–134, 1991.<br />
# I. Lenz, R. Knepper, and A. Saxena. Deepmpc: Learning deep latent features for model predictive control. In Robotics Science and Systems (RSS), 2015.<br />
# G. Wayne and L. Abbott. Hierarchical control using networks trained with higher-level forward models. Neural computation, 2014.<br />
# B. C. Stadie, S. Levine, and P. Abbeel. Incentivizing exploration in reinforcement learning with deep predictive models. arXiv preprint arXiv:1507.00814, 2015.<br />
# J. Clarke, D. Goldwasser, M.-W. Chang, and D. Roth. Driving semantic parsing from the world’s response. In Proceedings of the fourteenth conference on computational natural language learning, pages 18–27. Association for Computational Linguistics, 2010.<br />
# J. Schatzmann, K. Weilhammer, M. Stuttle, and S. Young. A survey of statistical user simulation techniques for reinforcement-learning of dialogue management strategies. The knowledge engineering review, 21(02):97–126, 2006.</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=Dialog-based_Language_Learning&diff=29759Dialog-based Language Learning2017-11-09T00:40:16Z<p>Nkadiyar: /* Learning models */</p>
<hr />
<div>This page is a summary for NIPS 2016 paper <i>Dialog-based Language Learning</i> [1].<br />
==Introduction==<br />
One of the ways humans learn language, especially second language or language learning by students, is by communication and getting its feedback. However, most existing research in Natural Language Understanding has focused on supervised learning from fixed training sets of labeled data. This kind of supervision is not realistic of how humans learn, where language is both learned by, and used for, communication. When humans act in dialogs (i.e., make speech utterances) the feedback from other human’s responses contain very rich information. This is perhaps most pronounced in a student/teacher scenario where the teacher provides positive feedback for successful communication and corrections for unsuccessful ones. <br />
<br />
This paper is about dialog-based language learning, where supervision is given naturally and implicitly in the response of the dialog partner during the conversation. This paper is a step towards the ultimate goal of being able to develop an intelligent dialog agent that can learn while conducting conversations. Specifically this paper explores whether we can train machine learning models to learn from dialog.<br />
<br />
===Contributions of this paper===<br />
*Introduce a set of tasks that model natural feedback from a teacher and hence assess the feasibility of dialog-based language learning. <br />
*Evaluated some baseline models on this data and compared them to standard supervised learning. <br />
*Introduced a novel forward prediction model, whereby the learner tries to predict the teacher’s replies to its actions, which yields promising results, even with no reward signal at all<br />
<br />
==Background on Memory Networks==<br />
<br/><br />
[[File:DB F1.png|center|700px]]<br />
<br/><br />
A memory network combines learning strategies from the machine learning literature with a memory component that can be read and written to.<br />
<br />
The high-level view of a memory network is as follows:<br />
*There is a memory, $m$, an indexed array of objects (e.g. vectors or arrays of strings).<br />
*An input feature map $I$, which converts the incoming input to the internal feature representation<br />
*A generalization component $G$ which updates old memories given the new input. <br />
*An output feature map $O$, which produces a new output in the feature representation space given the new input and the current memory state.<br />
*A response component $R$ which converts the output into the response format desired – for example, a textual response or an action.<br />
<br />
$I$, $G$, $O$ and $R$ can all potentially be learned components and make use of any ideas from the existing machine learning literature.<br />
<br />
In question answering systems for example, the components may be instantiated as follows:<br />
*$I$ can make use of standard pre-processing such as parsing, coreference, and entity resolution. It could also encode the input into an internal feature representation by converting from text to a sparse or dense feature vector.<br />
*The simplest form of $G$ is to introduce a function $H$ which maps the internal feature representation produced by I to an individual memory slot, and just updates the memory at $H(I(x))$.<br />
*$O$ Reads from memory and performs inference to deduce the set of relevant memories needed to perform a good response.<br />
*$R$ would produce the actual wording of the question answer based on the memories found by $O$. For example, $R$ could be an RNN conditioned on the output of $O$<br />
<br />
When the components $I$,$G$,$O$, & $R$ are neural networks, the authors describe the resulting system as a <b>Memory Neural Network (MemNN)</b>. They build a MemNN for QA (question answering) problems and compare it to RNNs (Recurrent Neural Network) and LSTMs (Long Short Term Memory RNNs) and find that it gives superior performance.<br />
<br />
[[File:DB_F2.png|center|800px]]<br />
<br />
==Related Work==<br />
<br />
'''Usefulness of feedback in language learning:''' Social interaction and natural infant directed conversations are shown to useful for language learning[2]. Several studies[3][4][5][6] have shown that feedback is especially useful in second language learning and learning by students.<br />
<br />
'''Supervised learning from dialogs using neural models:''' Neural networks has been used for response generation that can be trained end to end on large quantities of unstructured Twitter conversations[7]. However this does not incorporate feedback from dialog partner during real time conversation<br />
<br />
'''Reinforcement learning:''' Reinforcement learning works on dialogs[8][9], often consider reward as the feedback model rather than exploiting the dialog feedback per se.<br />
<br />
'''Forward prediction models:''' Although forward prediction models, have been used in other applications like learning eye-tracking[10], controlling robot arms[11] and vehicles[12], it has not been used for dialog.<br />
<br />
==Dialog-based Supervision tasks==<br />
For testing their models, the authors chose two datasets (i) the single supporting fact problem from the bAbI datasets [13] which consists of short stories from a simulated world followed by questions; and (ii) the MovieQA dataset [14] which is a large-scale dataset (∼ 100k questions over ∼ 75k entities) based on questions with answers in the open movie database (OMDb)<br />
<br />
However, since these datasets were not designed to model the supervision from dialogs, the authors modified them to create 10 supervision task types on these datasets(Fig 3).<br />
<br />
[[File:DB F3.png|center|700px]]<br />
<br/><br />
<br />
*'''Task 1: Imitating an Expert Student''': The dialogs take place between a teacher and an expert student who gives semantically coherent answers. Hence, the task is for the learner to imitate that expert student, and become an expert themselves <br />
<br />
*'''Task 2: Positive and Negative Feedback:''' When the learner answers a question the teacher then replies with either positive or negative feedback. In the experiments the subsequent responses are variants of “No, that’s incorrect” or “Yes, that’s right”. In the datasets there are 6 templates for positive feedback and 6 templates for negative feedback, e.g. ”Sorry, that’s not it.”, ”Wrong”, etc. To distinguish the notion of positive from negative, an additional external reward signal that is not part of the text<br />
<br />
*'''Task 3: Answers Supplied by Teacher:''' The teacher gives positive and negative feedback as in Task 2, however when the learner’s answer is incorrect, the teacher also responds with the correction. For example if “where is Mary?” is answered with the incorrect answer “bedroom” the teacher responds “No, the answer is kitchen”’<br />
<br />
*'''Task 4: Hints Supplied by Teacher:''' The corrections provided by the teacher do not provide the exact answer as in Task 3, but only a useful hint. This setting is meant to mimic the real life occurrence of being provided only partial information about what you did wrong.<br />
<br />
*'''Task 5: Supporting Facts Supplied by Teacher:''' Another way of providing partial supervision for an incorrect answer is explored. Here, the teacher gives a reason (explanation) why the answer is wrong by referring to a known fact that supports the true answer that the incorrect answer may contradict. <br />
<br />
*'''Task 6: Partial Feedback:''' External rewards are only given some of (50% of) the time for correct answers, the setting is otherwise identical to Task 3. This attempts to mimic the realistic situation of some learning being more closely supervised (a teacher rewarding you for getting some answers right) whereas other dialogs have less supervision (no external rewards). The task attempts to assess the impact of such partial supervision.<br />
<br />
*'''Task 7: No Feedback:''' External rewards are not given at all, only text, but is otherwise identical to Tasks 3 and 6. This task explores whether it is actually possible to learn how to answer at all in such a setting.<br />
<br />
*'''Task 8: Imitation and Feedback Mixture:''' Combines Tasks 1 and 2. The goal is to see if a learner can learn successfully from both forms of supervision at once. This mimics a child both observing pairs of experts talking (Task 1) while also trying to talk (Task 2).<br />
<br />
*'''Task 9: Asking For Corrections:''' The learner will ask questions to the teacher about what it has done wrong. Task 9 tests one of the most simple instances, where asking “Can you help me?” when wrong obtains from the teacher the correct answer.<br />
<br />
*'''Task 10: Asking for Supporting Facts:''' A less direct form of supervision for the learner after asking for help is to receive a hint rather than the correct answer, such as “A relevant fact is John moved to the bathroom” when asking “Can you help me?”. This is thus related to the supervision in Task 5 except the learner must request help<br />
<br />
[[File:F4.png|center|700px]]<br />
<br/><br />
<br />
For each task a fixed policy is considered for performing actions (answering questions) which gets questions correct with probability $π_{acc}$ (i.e. the chance of getting the red text correct in Figs. 3 and 4). We thus can compare different learning algorithms for each task over different values of $π_{acc}$ (0.5, 0.1 and 0.01). In all cases a training, validation and test set is provided.<br />
<br />
==Learning models==<br />
This work evaluates four possible learning strategies for each of the 10 tasks: imitation learning, reward-based imitation, forward prediction, and a combination of reward-based imitation and forward prediction<br />
<br />
All of these approaches are evaluated with the same model architecture: an end-to-end memory network (MemN2N) [15], which has been used as a baseline model for exploring differing modes of learning.<br />
<br />
[[File:F5.png|center|700px]]<br />
<br/><br />
The input is the last utterance of the dialog, $x$, as well as a set of memories (context) (<math> c_1</math>, . . . , <math> c_n</math> ) which can encode both short-term memory, e.g. recent previous utterances and replies, and long-term memories, e.g. facts that could be useful for answering questions. The context inputs <math> c_i</math> are converted into vectors <math> m_i</math> via embeddings and are stored in the memory. The goal is to produce an output $\hat{a}$ by processing the input $x$ and using that to address and read from the memory, $m$, possibly multiple times, in order to form a coherent reply. In the figure the memory is read twice, which is termed multiple “hops” of attention. <br />
<br />
In the first step, the input $x$ is embedded using a matrix $A$ of size $d$ × $V$ where $d$ is the embedding dimension and $V$ is the size of the vocabulary, giving $q$ = $A$$x$, where the input $x$ is as a bag-of words vector. Each memory <math> c_i</math> is embedded using the same matrix, giving $m_i$ = $A$$c_i$ . The output of addressing and then reading from memory in the first hop is: <br />
<br />
[[File:eq1.png|center|400px]]<br />
<br />
Here, $p^{1}$ is a probability vector over the memories, and is a measure of how much the input and the memories match. The goal is to select memories relevant to the last utterance $x$, i.e. the most relevant have large values of $p^{1}_i$ . The output memory representation $o_1$ is then constructed using the weighted sum of memories, i.e. weighted by $p^{1}$ . The memory output is then added to the original input, <math> c_1</math> = <math> R_1</math>(<math> o_1</math> + $q$), to form the new state of the controller, where <math> R_1</math> is a $d$ × $d$ rotation matrix . The attention over the memory can then be repeated using <math> u_1</math> instead of $q$ as the addressing vector, yielding: <br />
<br />
[[File:eq2.png|center|400px]]<br />
<br />
The controller state is updated again with <math> u_2</math> = <math> R_2</math>(<math> o_2</math> + <math> u_1</math>), where <math> R_2</math> is another $d$ × $d$ matrix to be learnt. In a two-hop model the final output is then defined as: <br />
<br />
[[File:eq3.png|center|400px]]<br />
<br />
where there are $C$ candidate answers in $y$. In our experiments $C$ is the set of actions that occur in the training set for the bAbI tasks, and for MovieQA it is the set of words retrieved from the KB.<br />
<br />
==Training strategies==<br />
*'''1. Imitation Learning'''<br />
This approach involves simply imitating one of the speakers in observed dialogs. Examples arrive as $(x, c, a)$ triples, where $a$ is (assumed to be) a good response to the last utterance $x$ given context $c$. Here, the whole memory network model defined above is trained using stochastic gradient descent by minimizing a standard cross-entropy loss between $a^{ˆ}$ and the label $a$<br />
<br />
*'''2. Reward-based Imitation''' <br />
If some actions are poor choices, then one does not want to repeat them, that is we shouldn’t treat them as a supervised objective. Here, positive reward is only obtained immediately after (some of) the correct actions, or else is zero. Only apply imitation learning on the rewarded actions. The rest of the actions are simply discarded from the training set. <br />
<br />
*'''3. Forward Prediction''' <br />
The aim is, given an utterance $x$ from speaker 1 and an answer a by speaker 2 (i.e., the learner), to predict $x^{¯}$, the response to the answer from speaker 1. That is, in general to predict the changed state of the world after action $a$, which in this case involves the new utterance $x^{¯}$.<br />
<br />
[[File:F6.png|center|700px]]<br />
<br/><br />
As shown in Figure (b), this is achieved by chopping off the final output from the original network of Fig (a) and replace it with some additional layers that compute the forward prediction. The first part of the network remains exactly the same and only has access to input x and context c, just as before. The computation up to $u_2$ = $R_2$($o_2$ + $u_1$) is thus exactly the same as before. <br />
<br />
Then perform another “hop” of attention but over the candidate answers rather than the memories. The information of which action (candidate) was actually selected in the dialog (i.e. which one is a) is also incorporated which is crucial. After this “hop”, the resulting state of the controller is then used to do the forward prediction.<br />
<br />
Concretely, we compute: <br />
<br />
[[File:eq4.png|center|550px]]<br />
<br />
where $β^{*}$ is a d-dimensional vector, that is also learnt, that represents in the output $o_3$ the action that was actually selected. After obtaining $o_3$, the forward prediction is then computed as: <br />
<br />
[[File:eq5.png|center|500px]]<br />
<br />
where $u_3$ = $R_3$($o_3$ + $u_2$). That is, it computes the scores of the possible responses to the answer a over $C^{¯}$ possible candidates.<br />
<br />
<br />
*'''4. Reward-based Imitation + Forward Prediction'''<br />
As the reward-based imitation learning uses the architecture of Fig (a), and forward prediction uses the same architecture but with the additional layers of Fig (b), we can learn jointly with both strategies. This is a powerful combination as it makes use of reward signal when available and the dialog feedback when the reward signal is not available.<br />
<br />
==Experiments==<br />
<br />
Experiments were conducted on the two test datasets - babi and MovieQA. For each task, a fixed policy is considered for performing actions (answering questions) which gets questions correct with probability $π_{acc}$. This helps to compare the different training strategies described earlier over each task for different values of $π_{acc}$. Hyperparameters for all methods are optimized on the validation sets.<br />
<br/><br/><br />
[[File:DB F7.png|center|800px]]<br />
<br/><br />
The following results are observed by the authors:<br />
*Imitation learning, ignoring rewards, is a poor learning strategy when imitating inaccurate answers, e.g. for $π_{acc}$ < 0.5. For imitating an expert however (Task 1) it is hard to beat. <br />
*Reward-based imitation (RBI) performs better when rewards are available, particularly in Table 1, but also degrades when they are too sparse e.g. for πacc = 0.01.<br />
*Forward prediction (FP) is more robust and has stable performance at different levels of πacc. However as it only predicts answers implicitly and does not make use of rewards it is outperformed by RBI on several tasks, notably Tasks 1 and 8 (because it cannot do supervised learning) and Task 2 (because it does not take advantage of positive rewards).<br />
*FP makes use of dialog feedback in Tasks 3-5 whereas RBI does not. This explains why FP does better with useful feedback (Tasks 3-5) than without (Task 2), whereas RBI cannot.<br />
*Supplying full answers (Task 3) is more useful than hints (Task 4) but hints still help FP more than just yes/no answers without extra information (Task 2).<br />
*When positive feedback is sometimes missing (Task 6) RBI suffers especially in Table 1. FP does not as it does not use this feedback.<br />
*One of the most surprising results of our experiments is that FP performs well overall, given that it does not use feedback, which we will attempt to explain subsequently. This is particularly evident on Task 7 (no feedback) where RBI has no hope of succeeding as it has no positive examples. FP on the other hand learns adequately.<br />
*Tasks 9 and 10 are harder for FP as the question is not immediately before the feedback.<br />
*Combining RBI and FP ameliorates the failings of each, yielding the best overall results<br />
<br />
==Future work==<br />
* Any reply in a dialog can be seen as feedback, and should be useful for learning. Evaluate if forward prediction, and the other approaches in this paper, work there too. <br />
* Develop further evaluation methodologies to test how the models presented here work in more complex settings where actions that are made lead to long-term changes in the environment and delayed rewards, i.e. extending to the reinforcement learning setting, and to full language generation. <br />
* How dialog-based feedback could also be used as a medium to learn non-dialog based skills, e.g. natural language dialog for completing visual or physical tasks.<br />
<br />
==References==<br />
# Jason Weston. Dialog-based Language Learning. NIPS, 2016.<br />
# P. K. Kuhl. Early language acquisition: cracking the speech code. Nature reviews neuroscience, 5(11): 831–843, 2004.<br />
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# J. Schatzmann, K. Weilhammer, M. Stuttle, and S. Young. A survey of statistical user simulation techniques for reinforcement-learning of dialogue management strategies. The knowledge engineering review, 21(02):97–126, 2006.</div>Nkadiyarhttp://wiki.math.uwaterloo.ca/statwiki/index.php?title=STAT946F17/Cognitive_Psychology_For_Deep_Neural_Networks:_A_Shape_Bias_Case_Study&diff=29757STAT946F17/Cognitive Psychology For Deep Neural Networks: A Shape Bias Case Study2017-11-09T00:30:37Z<p>Nkadiyar: /* Introduction */</p>
<hr />
<div>= Introduction =<br />
<br />
The recent burgeon on the use of Deep Neural Networks (DNNs) have resulted in giant leaps of accuracy in prediction. They are also being used to solve a variety of complex tasks which earlier methodologies have struggled to excel in.<br />
<br />
While it is all good to see incredibly high accuracy as a result of the use of DNN, we must begin to question why they perform so well. It has become an interesting field of study to actually represent the features/feature maps or interpret the meaning of the learnt values in a DNN's hidden layers. Currently we treat models of DNNs as black boxes which we practically tune the tweakable parameters like number of layers, number of units in each layer, number & size of feature maps(in case of CNN) etc. The opacity created by the lack of an intuitive representation of the internal learnt parameters of DNNs hinders both basic research as well as its application to real world problems.<br />
<br />
Recent pushes have aimed to better understand DNNs: tailor-made loss functions and architectures produce more interpretable features (Higgins et al., 2016; Raposo et al., 2017) while output-behavior analyses unveil previously opaque operations of these networks (Karpathy et al., 2015). Parallel to this work, neuroscience-inspired methods such as activation visualization (Li et al., 2015), ablation analysis (Zeiler & Fergus, 2014) and activation maximization (Yosinski et al., 2015) have also been applied<br />
<br />
This paper aims to provide another methodology to attempt to decipher & better understand how DNNs solve a particular task. This methodology was inspired by psychological concepts to test whether the DNN's were able to make accurate predictions with biases similar to that the human mind makes.<br />
<br />
Research in developmental psychology shows that when learning new words, humans tend to assign the same name to similarly shaped items rather than to items with similar color, texture, or size. This bias/knowledge tend to be forged into the brains of humans and humans then take this forward to easily associate these shapes with new objects they have not seen before.<br />
<br />
The authors of this paper try to test if DNNs behave similarly in one-shot learning applications. They attempt to prove that when the models of state-of-the-art DNNs are used to learn objects from images, they exhibit a stronger shape bias than a color bias. To emulate the human brain, they use the parameters of pre-trained DNN models and use this to perform one-shot learning on a new data set with different labels.<br />
<br />
= Background =<br />
== One Shot Learning ==<br />
One-shot learning is an object categorization problem in computer vision. Whereas most machine learning based object categorization algorithms require training on hundreds or thousands of images and very large datasets, one-shot learning aims to learn information about object categories from one, or only a few, training images.<br />
<br />
The one-shot word learning task is to label a novel data example $\hat{x}$ (e.g. a novel probe image) with a novel class label $\hat{y}$ (e.g. a new word) after only a single example.<br />
<br />
More specifically, given a support set $S = {(x_i, y_i) , i \in [1, k]}$, of images $x_i$, and their associated labels $y_i$, and an unlabeled probe image $\hat{x}$,<br />
the one-shot learning task is to identify the true label of the probe image, $\hat{y}$, from the support set labels $ {y_i , i \in [1, k]} $:<br />
<br />
<br />
$\displaystyle \hat{y} = arg \max_{y}$ $P(y | \hat{x}, S)$<br />
<br />
<br />
We assume that the image labels $y_i$ are represented using a one-hot encoding and that $P(y|\hat{x}, S)$ is parameterised by a DNN, allowing us to leverage the ability of deep networks to learn powerful representations.<br />
<br />
== Inception Networks ==<br />
<br />
A probe image $\hat{x}$ is given the label of the nearest neighbour from the<br />
support set:<br />
<br />
$\hat{y} = y$<br />
<br />
$(x, y) = \displaystyle arg \min_{(x_i,y_i) \in S} d(h(x_i), h(\hat{x})) $<br />
<br />
where d is a distance function.<br />
<br />
The function h is parameterized by Inception – one of the best performing ImageNet classification models. Specifically, h returns features from the last layer (the softmax input) of a pre-trained Inception classifier. With these features as input and cosine distance as the distance function, the classifier in achieves 87.6% accuracy on one-shot classification on the ImageNet dataset (Vinyals et al., 2016). We call the Inception classifier together with the nearest-neighbor component the Inception Baseline (IB) model.<br />
<br />
<br />
== Matching Networks ==<br />
<br />
MNs (Vinyals et al.,2016) are neural network architectures with state-of-the-art one shot learning performance on ImageNet (93.2% one-shot labelling accuracy).<br />
MNs are trained to assign label $\hat{y}$ to probe image $\hat{x}$ using an attention mechanism a acting on image embeddings stored in the support set S:<br />
\begin{align*}<br />
a(\hat{x},x_i)=\frac{e^{d(f(\hat{x},S), g(x_i,S))}}{\sum_{j}e^{d(f(\hat{x},S), g(x_j,S))}},<br />
\end{align*}<br />
<br />
<br />
where d is a cosine distance and where f and g provide context-dependent embeddings of $\hat{x}$ and $x_i$ (with contextS). The embedding $g(x_i, S)$ is a bi-directional LSTM (Hochreiter & Schmidhuber, 1997) with the support set S provided as an input sequence. The embedding $f(\hat{x}, S)$ is an LSTM with a read-attention mechanism operating over the entire embedded support set. The input to the LSTM is given by the penultimate layer features of a pre-trained deep convolutional network, specifically Inception.<br />
<br />
To train MNs we proceed as follows:<br />
<br />
=== Training MN ===<br />
* Step 1: At each step of training, the model is given a small support set of images and associated labels. In addition to the support set, the model is fed an unlabeled probe image $\hat{x}$<br />
<br />
* Step 2: The model parameters are then updated to improve classification accuracy of the probe image $\hat{x}$ given the support set. Parameters are updated using stochastic gradient descent with a learning rate of 0.1<br />
<br />
* Step 3: After each update, the labels ${(y_i, i \in [1, k]}$ in the training set are randomly re-assigned to new image classes (the label indices are randomly permuted,<br />
but the image labels are not changed). This is a critical step. It prevents MNs from learning a consistent mapping between a category and a label. Usually, in classification, this is what we want, but in one-shot learning we want to train our model for classification after viewing a single in-class example from the support set. <br />
<br />
The objective function used is:<br />
<br />
\begin{align*}<br />
L=E_{C\sim T}\biggr[E_{S \sim C, B \sim C}\bigr[\sum_{(x,y)\in B}\log P(y|x,S)\bigr]\biggr]<br />
\end{align*}<br />
<br />
<br />
where T is the set of all possible labelings of our classes, S is a support set sampled with a class labeling C ~ T and B is a batch of probe images and labels, also with the same randomly chosen class labeling as the support set.<br />
<br />
== Cognitive Biases ==<br />
Cognitive bias is a concept from developmental psychology which attempts to explain how children can extract meanings of words with very few examples, similar to the concept of one-shot learning discussed above. The theory, as explained by the authors, is that humans form biases that allow them to eliminate many potential hypotheses about word meaning where the amount of data available is insufficient for this purpose. These include:<br />
* Whole object bias<br />
* Taxonomic bias<br />
* Mutual exclusivity bias<br />
* Shape bias<br />
A more complete list of cognitive biases is given by [[#References|(Bloom, 2000)]]. The bias the authors investigate in this paper is the shape bias, which denotes a tendency to assign the same name to similarly shaped items rather than to items with similar color, texture, or size.<br />
<br />
= Methodology =<br />
== Inductive Biases & Probe Data ==<br />
<br />
Inductive biases are those criteria which are artificially selected or learnt by the network as a classifying/distinguishing property.<br />
It has been observed that the biases that DNNs learnt are complex composite features. We, as researchers can take advantage of the fact that DNNs learnt complex distinguishing features by constructing probe data sets which particularly target on exposing a particular bias that a DNN might have. <br />
<br />
* Step 1: Take a known composite feature which we suspect the DNNs are biased against. It is possible that the feature is not numerical but intuitive for human researchers to understand.<br />
* Step 2: Train the target model with an appropriate dataset.<br />
* Step 3: Transfer Learning: Use the pre-trained model with a new data set which is curated to contain data to prove/disprove the existence of the bias<br />
* Step 4: Model/Decide on a function which quantifies the bias under study<br />
* Step 5: Measure the bias with the bias function<br />
<br />
Then it is possible to deduce whether the DNN uses the feature to solve the task by the value of bias function.<br />
<br />
== Data Sets Used ==<br />
<br />
* Training Set: ImageNet<br />
* Test Set:<br />
** The Cognitive Psychology Probe Data (CogPsyc data) that is used consists of 150 images of objects. The images are arranged in triples consisting of a probe image, a shape-match image (that matches the probe in colour but not shape), and a color-match image (that matches the probe in shape but not colour). In the dataset there are 10 triples, each shown on 5 different backgrounds, giving a total of 50 triples. [[File:CogPsy.PNG|center|350px]]<br />
** A real-world dataset consisting of 90 images of objects (30 triples) collected using Google Image Search. The images are arranged in triples consisting of a probe, a shape-match and a colour-match.<br />
<br />
= Experiments =<br />
== Evaluation Criteria ==<br />
<br />
* For a given probe image $\hat{x}$, we loaded the shape-match image $x_s$ and corresponding label $y_s$, along with the colour-match image $x_c$ and corresponding label $y_c$ into memory, as the support set $S = \{(x_s, ys), (x_c, y_c)\}$<br />
* Calculate $\hat{y}$<br />
* The model assigns either $y_c$ or $y_s$ to the probe image.<br />
* To estimate the shape bias Bs, calculate the proportion of shape labels assigned to the probe: $B_s = E(\delta(\hat{y} - y_s))$ <br />
where E is an expectation across probe images and $\delta$ is the Dirac delta function.<br />
<br />
== Experiment 1: Shape bias statistics in Inception Baseline: ==<br />
* Shape bias of IB to be $B_s = 0.68$. Similarly, the shape bias of IB using our real-world dataset was $B_s = 0.97$. Together, these results strongly suggest that IB trained on ImageNet has a stronger bias towards shape than colour<br />
<br />
== Experiment 2: Shape bias statistics in Matching Network: ==<br />
* They found that MNs have a shape of bias $B_s = 0.7$ using the CogPsyc dataset and a bias of $Bs = 1$ using the real-world dataset. Once again, these results suggest that MNs trained seeding from Inception using ImageNet has a stronger bias towards shape than colour.<br />
<br />
== Experiment 3: Shape bias statistics between and across models: ==<br />
<br />
The authors extended the shape bias analysis to calculate the shape bias in a population of IB models and in a population of MN models with different random initialization<br />
<br />
=== Dependence on the initialization of parameters: ===<br />
<br />
[[File:3.1.PNG|right|250px]] A strong variability was observed when variation in the initial values of the parameters. For the CogPsyc dataset, the average shape bias was $B_s = 0.628$ with standard deviation $\sigma B_s = 0.049$ at the end of training and for the real-world dataset the average shape bias was $B_s = 0:958$ with $\sigma B_s = 0.037$.<br />
<br />
=== Dependence of shape bias on model performance: ===<br />
<br />
For the CogPsych dataset, the correlation between bias and classification accuracy was $\rho = 0.15$, and for the real world dataset, correlation between bias and classification accuracy was $\rho = -0.06$. This would be evident since the accuracy of the models remained nearly constant when the initialization parameters varied whereas the shape bias tended to vary a lot, hence highlighting the lack of correlation amongst them. <br />
<br />
=== Emergence of shape bias during training: ===<br />
The shape bias spiked to a large value very early. <br />
<br />
=== Variation of shape bias within models & across models: ===<br />
With different initialization parameters, the shape bias varied a lot within IB during training while the shape bias did not fluctuate during the training of MN. It was found that the MN inherits the shape bias of the IB which seeded its embeddings and thereafter, the shape bias remained constant throughout training. It is important to note that the output of the penultimate layer of the Inception was not fine tuned when it was pipelined to the MN. This was to ensure that the MN properties were independent of the IB model properties. [[File:3.3.PNG|center|250px]] [[File:3.4.PNG|center|250px]]<br />
<br />
= Learnings, Inferences & Implications =<br />
* Both the Inception Baseline and the Matching Network exhibit strong shape bias when trained on ImageNet. Researchers who use Inception & MN DNNs can now use this fact as a consideration for their application while using pre-trained models for new datasets. If it is known before hand that the new data set is strongly classifiable through a color bias, then they would either want to defer using the pre-trained models or explore methods to decrease/remove the strong shape bias.<br />
<br />
* There exists a high variability in the shape bias with the variation in the initialization parameters. This is an important finding since it uncovers the fact that the same architecture which exhibit similar accuracy in predictions can display a variety of shape bias just with different initialization parameters. Researchers can explore methods of tuning the random initialization such that the models start out with a low shape bias without compromising the accuracy of the model.<br />
<br />
* MNs inherit the shape bias which is seeded to it by the Inception Network's input embedding. This is also another fact which researchers & practitioners should be careful about. When using cascaded or pipelined heterogeneous architectures, the models downstream tend to inherit/become/are fed with the properties/biases of the models upstream. This may be desirable or undesirable according to the application, but it is important to be aware of its presence.<br />
<br />
* The biases under consideration are the property of the collection of the architecture, the dataset and the optimization procedure. Hence in order to increase or decrease the effect of a particular bias, one or more of the mentioned factors must be adjusted/tuned/changed.<br />
<br />
* The fact that a high shape bias emerged in the early epochs with less variability in further epochs can be thought of analogous to the biases that humans develop at an infancy which gets fortified as they age.<br />
<br />
=== Modeling human word learning === <br />
The authors note how there have been several previous attempts to model human word learning in the field of cognitive science. A major shortcoming of these works is that none of the models are capable of one-shot word learning when presented with real-world images. Recognizing this and given the success of MNs in this work, the authors propose MNs as a computational-level account of one-shot word learning. <br />
In the paper, it is discussed how shape bias increases dramatically during early training of the model. This agrees with what is observed in many psychological studies: older children show bias more than younger children. <br />
<br />
= Conclusion, Future Work and Open questions =<br />
<br />
* Just as cognitive psychology exposes the shape bias observed in this experiment, we should try to uncover other biases as well using multiple approaches<br />
* Study the underlying mechanisms which cause biases such as shape bias in DNNs<br />
* Research into various methods of probing and creating probe data sets which can be used to test architectures for various biases<br />
* Exploration into a research field called Artificial Cognitive Psychology which focuses on probing how DNN architectures can be understood further using known behaviors of the human brain<br />
<br />
= References =<br />
<br />
* Ritter, Samuel & G. T. Barrett, David & Santoro, Adam & M. Botvinick, Matt. (2017). Cognitive Psycholo